* add: a workflow to check if docker containers can be built if the files are modified.
* type
* unify docker image build test and push
* make it run on prs too.
* check
* check
* check
* check again.
* remove docker test build file.
* remove extra dependencies./
* check
* Initial commit
* Removed copy hints, as in original SDXLControlNetPipeline
Removed copy hints, as in original SDXLControlNetPipeline, as the `make fix-copies` seems to have issues with the @property decorator.
* Reverted changes to ControlNetXS
* Addendum to: Removed changes to ControlNetXS
* Added test+docs for mixture of denoiser
* Update docs/source/en/using-diffusers/controlnet.md
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update docs/source/en/using-diffusers/controlnet.md
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
* Update unet_2d_blocks.py
similar doc-string add to have in the original diffusion repository.
* Fix bug for mention in this issue section #6901
* Update src/diffusers/schedulers/scheduling_ddim_flax.py
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* Fix linter
* Restore empty line
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* copied from for t2i pipelines without ip adapter support.
* two more pipelines with proper copied from comments.
* revert to the original implementation
* throw error when patch inputs and layernorm are provided for transformers2d.
* add comment on supported norm_types in transformers2d
* more check
* fix: norm _type handling
* [bug] Fix float/int guidance scale not working in `StableVideoDiffusionPipeline`
* Add test to disable CFG on SVD
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* support and example launch for sdxl turbo
* White space fixes
* Trailing whitespace character
* ruff format
* fix guidance_scale and steps for turbo mode
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Radames Ajna <radamajna@gmail.com>
* update svd docs
* fix example doc string
* update return type hints/docs
* update type hints
* Fix typos in pipeline_stable_video_diffusion.py
* make style && make fix-copies
* Update src/diffusers/pipelines/stable_video_diffusion/pipeline_stable_video_diffusion.py
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update src/diffusers/pipelines/stable_video_diffusion/pipeline_stable_video_diffusion.py
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* update based on suggestion
---------
Co-authored-by: M. Tolga Cangöz <mtcangoz@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Enable FakeTensorMode for EulerDiscreteScheduler scheduler
PyTorch's FakeTensorMode does not support `.numpy()` or `numpy.array()`
calls.
This PR replaces `sigmas` numpy tensor by a PyTorch tensor equivalent
Repro
```python
with torch._subclasses.FakeTensorMode() as fake_mode, ONNXTorchPatcher():
fake_model = DiffusionPipeline.from_pretrained(model_name, low_cpu_mem_usage=False)
```
that otherwise would fail with
`RuntimeError: .numpy() is not supported for tensor subclasses.`
* Address comments
* add tags for diffusers training
* add tags for diffusers training
* add tags for diffusers training
* add tags for diffusers training
* add tags for diffusers training
* add tags for diffusers training
* add dora tags for drambooth lora scripts
* style
* add is_dora arg
* style
* add dora training feature to sd 1.5 script
* added notes about DoRA training
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* initial
* check_inputs fix to the rest of pipelines
* add fix for no cfg too
* use of variable
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Add copyright notice to relevant files and fix typos
* Set `timestep_spacing` parameter of `StableDiffusionXLPipeline`'s scheduler to `'trailing'`.
* Update `StableDiffusionXLPipeline.from_single_file` by including EulerAncestralDiscreteScheduler with `timestep_spacing="trailing"` param.
* Update model loading method in SDXL Turbo documentation
* move model helper function in pipeline to EfficiencyMixin
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* DPMMultistep rescale_betas_zero_snr
* DPM upcast samples in step()
* DPM rescale_betas_zero_snr UT
* DPMSolverMulti move sample upcast after model convert
Avoids having to re-use the dtype.
* Add a newline for Ruff
* log_validation unification for controlnet.
* additional fixes.
* remove print.
* better reuse and loading
* make final inference run conditional.
* Update examples/controlnet/README_sdxl.md
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* resize the control image in the snippet.
---------
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* Make LoRACompatibleConv padding_mode work.
* Format code style.
* add fast test
* Update src/diffusers/models/lora.py
Simplify the code by patrickvonplaten.
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* code refactor
* apply patrickvonplaten suggestion to simplify the code.
* rm test_lora_layers_old_backend.py and add test case in test_lora_layers_peft.py
* update test case.
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* modulize log validation
* run make style and refactor wanddb support
* remove redundant initialization
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* make checkpoint_merger pipeline pass the "variant" argument to from_pretrained()
* make style
---------
Co-authored-by: Lincoln Stein <lstein@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* add stable_diffusion_xl_ipex community pipeline
* make style for code quality check
* update docs as suggested
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* standardize model card
* fix tags
* correct import styling and update tags
* run make style and make quality
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* feat: allow low_cpu_mem_usage in ip adapter loading
* reduce the number of device placements.
* documentation.
* throw low_cpu_mem_usage warning only once from the main entry point.
* use load_model_into_meta in single file utils
* propagate to autoencoder and controlnet.
* correct class name access behaviour.
* remove torch_dtype from load_model_into_meta; seems unncessary
* remove incorrect kwarg
* style to avoid extra unnecessary line breaks
* fix: bias loading bug
* fixes for SDXL
* apply changes to the conversion script to match single_file_utils.py
* do transpose to match the single file loading logic.
Remove <cat-toy> validation prompt from textual_inversion_sdxl.py
The `<cat-toy>` validation prompt is a default choice for the example task in the README. But no other part of `textual_inversion_sdxl.py` references the cat toy and `textual_inversion.py` has a default validation prompt of `None` as well.
So bring `textual_inversion_sdxl.py` in line with `textual_inversion.py` and change default validation prompt to `None`
* attention_head_dim
* debug
* print more info
* correct num_attention_heads behaviour
* down_block_num_attention_heads -> num_attention_heads.
* correct the image link in doc.
* add: deprecation for num_attention_head
* fix: test argument to use attention_head_dim
* more fixes.
* quality
* address comments.
* remove depcrecation.
* add: support for passing ip adapter image embeddings
* debugging
* make feature_extractor unloading conditioned on safety_checker
* better condition
* type annotation
* index to look into value slices
* more debugging
* debugging
* serialize embeddings dict
* better conditioning
* remove unnecessary prints.
* Update src/diffusers/loaders/ip_adapter.py
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* make fix-copies and styling.
* styling and further copy fixing.
* fix: check_inputs call in controlnet sdxl img2img pipeline
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* feat: standarize model card creation for dreambooth training.
* correct 'inference
* remove comments.
* take component out of kwargs
* style
* add: card template to have a leaner description.
* widget support.
* propagate changes to train_dreambooth_lora
* propagate changes to custom diffusion
* make widget properly type-annotated
* fix: callback function name is incorrect
On this tutorial there is a function defined and then used inside `callback_on_step_end` argument, but the name was not correct (mismatch)
* fix: typo in num_timestep (correct is num_timesteps)
fixed property name
* remove _to_tensor
* remove _to_tensor definition
* remove _collapse_frames_into_batch
* remove lora for not bloating the code.
* remove sample_size.
* simplify code a bit more
* ensure timesteps are always in tensor.
* Fix `AutoencoderTiny` with `use_slicing`
When using slicing with AutoencoderTiny, the encoder mistakenly encodes the entire batch for every image in the batch.
* Fixed formatting issue
* add noise_offset param
* micro conditioning - wip
* image processing adjusted and moved to support micro conditioning
* change time ids to be computed inside train loop
* change time ids to be computed inside train loop
* change time ids to be computed inside train loop
* time ids shape fix
* move token replacement of validation prompt to the same section of instance prompt and class prompt
* add offset noise to sd15 advanced script
* fix token loading during validation
* fix token loading during validation in sdxl script
* a little clean
* style
* a little clean
* style
* sdxl script - a little clean + minor path fix
sd 1.5 script - change default resolution value
* ad 1.5 script - minor path fix
* fix missing comma in code example in model card
* clean up commented lines
* style
* remove time ids computed outside training loop - no longer used now that we utilize micro-conditioning, as all time ids are now computed inside the training loop
* style
* [WIP] - added draft readme, building off of examples/dreambooth/README.md
* readme
* readme
* readme
* readme
* readme
* readme
* readme
* readme
* removed --crops_coords_top_left from CLI args
* style
* fix missing shape bug due to missing RGB if statement
* add blog mention at the start of the reamde as well
* Update examples/advanced_diffusion_training/README.md
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* change note to render nicely as well
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Fix bug in ResnetBlock2D.forward when not USE_PEFT_BACKEND and using scale_shift for time emb where the lora scale gets overwritten.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* fix minsnr implementation for v-prediction case
* format code
* always compute snr when snr_gamma is specified
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* feat: explicitly tag to diffusers when using push_to_hub
* remove tags.
* reset repo.
* Apply suggestions from code review
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix: tests
* fix: push_to_hub behaviour for tagging from save_pretrained
* Apply suggestions from code review
Co-authored-by: Lucain <lucainp@gmail.com>
* Apply suggestions from code review
Co-authored-by: Lucain <lucainp@gmail.com>
* import fixes.
* add library name to existing model card.
* add: standalone test for generate_model_card
* fix tests for standalone method
* moved library_name to a better place.
* merge create_model_card and generate_model_card.
* fix test
* address lucain's comments
* fix return identation
* Apply suggestions from code review
Co-authored-by: Lucain <lucainp@gmail.com>
* address further comments.
* Update src/diffusers/pipelines/pipeline_utils.py
Co-authored-by: Lucain <lucainp@gmail.com>
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Lucain <lucainp@gmail.com>
* initial commit for unconditional/class-conditional consistency training script
* make style
* Add entry for consistency training script in community README.
* Move consistency training script from community to research_projects/consistency_training
* Add requirements.txt and README to research_projects/consistency_training directory.
* Manually revert community README changes for consistency training.
* Fix path to script after moving script to research projects.
* Add option to load U-Net weights from pretrained model.
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* begin animatediff img2video and video2video
* revert animatediff to original implementation
* add img2video as pipeline
* update
* add vid2vid pipeline
* update imports
* update
* remove copied from line for check_inputs
* update
* update examples
* add multi-batch support
* fix __init__.py files
* move img2vid to community
* update community readme and examples
* fix
* make fix-copies
* add vid2vid batch params
* apply suggestions from review
Co-Authored-By: Dhruv Nair <dhruv.nair@gmail.com>
* add test for animatediff vid2vid
* torch.stack -> torch.cat
Co-Authored-By: Dhruv Nair <dhruv.nair@gmail.com>
* make style
* docs for vid2vid
* update
* fix prepare_latents
* fix docs
* remove img2vid
* update README to :main
* remove slow test
* refactor pipeline output
* update docs
* update docs
* merge community readme from :main
* final fix i promise
* add support for url in animatediff example
* update example
* update callbacks to latest implementation
* Update src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Update src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix merge
* Apply suggestions from code review
* remove callback and callback_steps as suggested in review
* Update tests/pipelines/animatediff/test_animatediff_video2video.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix import error caused due to unet refactor in #6630
* fix numpy import error after tensor2vid refactor in #6626
* make fix-copies
* fix numpy error
* fix progress bar test
---------
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* sd1.5 support in separate script
A quick adaptation to support people interested in using this method on 1.5 models.
* sd15 prompt text encoding and unet conversions
as per @linoytsaban 's recommendations. Testing would be appreciated,
* Readability and quality improvements
Removed some mentions of SDXL, and some arguments that don't apply to sd 1.5, and cleaned up some comments.
* make style/quality commands
* tracker rename and run-it doc
* Update examples/advanced_diffusion_training/train_dreambooth_lora_sd15_advanced.py
* Update examples/advanced_diffusion_training/train_dreambooth_lora_sd15_advanced.py
---------
Co-authored-by: Linoy Tsaban <57615435+linoytsaban@users.noreply.github.com>
* move unets to module 🦋
* parameterize unet-level import.
* fix flax unet2dcondition model import
* models __init__
* mildly depcrecating models.unet_2d_blocks in favor of models.unets.unet_2d_blocks.
* noqa
* correct depcrecation behaviour
* inherit from the actual classes.
* Empty-Commit
* backwards compatibility for unet_2d.py
* backward compatibility for unet_2d_condition
* bc for unet_1d
* bc for unet_1d_blocks
* Fixed the bug related to saving DeepSpeed models.
* Add information about training SD models using DeepSpeed to the README.
* Apply suggestions from code review
---------
Co-authored-by: mhh001 <mahonghao1@huawei.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* - extract function for stage in UNet2DConditionModel init & forward
- Add new function get_mid_block() to unet_2d_blocks.py
* add type hint to get_mid_block aligned with get_up_block and get_down_block; rename _set_xxx function
* add type hint and use keyword arguments
* remove `copy from` in versatile diffusion
* add animatediff img2vid
* fix
* Update examples/community/README.md
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix code snippet between ip adapter face id and animatediff img2vid
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* [Fix] Multiple image conditionings in a single batch for `StableDiffusionControlNetPipeline`.
* Refactor `check_inputs` in `StableDiffusionControlNetPipeline` to avoid redundant codes.
* Make the behavior of MultiControlNetModel to be the same to the original ControlNetModel
* Keep the code change minimum for nested list support
* Add fast test `test_inference_nested_image_input`
* Remove redundant check for nested image condition in `check_inputs`
Remove `len(image) == len(prompt)` check out of `check_image()`
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Better `ValueError` message for incompatible nested image list size
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Fix syntax error in `check_inputs`
* Remove warning message for multi-ControlNets with multiple prompts
* Fix a typo in test_controlnet.py
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Add test case for multiple prompts, single image conditioning in `StableDiffusionMultiControlNetPipelineFastTests`
* Improved `ValueError` message for nested `controlnet_conditioning_scale`
* Documenting the behavior of image list as `StableDiffusionControlNetPipeline` input
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Fixes#6418 Advanced Dreambooth LoRa Training
* change order of import to fix nit
* fix nit, use cast_training_params
* remove torch.compile fix, will move to a new PR
* remove unnecessary import
* Enable image resizing to adjust its height and width in StableDiffusionXLInstructPix2PixPipeline
* Ensure that validation is performed at every 'validation_step', not at every step
* fix: training resume from fp16.
* add: comment
* remove residue from another branch.
* remove more residues.
* thanks to Younes; no hacks.
* style.
* clean things a bit and modularize _set_state_dict_into_text_encoder
* add comment about the fix detailed.
* support compile
* make style
* move unwrap_model inside function
* change unwrap call
* run make style
* Update examples/dreambooth/train_dreambooth.py
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Revert "Update examples/dreambooth/train_dreambooth.py"
This reverts commit 70ab09732e.
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Remove conversion to RGB
* Add a Conversion Function
* Add type hint for convert_method
* Update src/diffusers/utils/loading_utils.py
Update docstring
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Update docstring
* Optimize imports
* Optimize imports (2)
* Reformat code
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* base template file - train_instruct_pix2pix.py
* additional import and parser argument requried for lora
* finetune only instructpix2pix model -- no need to include these layers
* inject lora layers
* freeze unet model -- only lora layers are trained
* training modifications to train only lora parameters
* store only lora parameters
* move train script to research project
* run quality and style code checks
* move train script to a new folder
* add README
* update README
* update references in README
---------
Co-authored-by: Rahul Raman <rahulraman@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* enable stable-xl textual inversion
* check if optimizer_2 exists
* check text_encoder_2 before using
* add textual inversion for sdxl in a single file
* fix style
* fix example style
* reset for error changes
* add readme for sdxl
* fix style
* disable autocast as it will cause cast error when weight_dtype=bf16
* fix spelling error
* fix style and readme and 8bit optimizer
* add README_sdxl.md link
* add tracker key on log_validation
* run style
* rm the second center crop
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* add tutorials to toctree.yml
* fix title
* fix words
* add overview ja
* fix diffusion to 拡散
* fix line 21
* add space
* delete supported pipline
* fix tutorial_overview.md
* fix space
* fix typo
* Delete docs/source/ja/tutorials/using_peft_for_inference.md
this file is not translated
* Delete docs/source/ja/tutorials/basic_training.md
this file is not translated
* Delete docs/source/ja/tutorials/autopipeline.md
this file is not translated
* fix toctree
* add: experimental script for diffusion dpo training.
* random_crop cli.
* fix: caption tokenization.
* fix: pixel_values index.
* fix: grad?
* debug
* fix: reduction.
* fixes in the loss calculation.
* style
* fix: unwrap call.
* fix: validation inference.
* add: initial sdxl script
* debug
* make sure images in the tuple are of same res
* fix model_max_length
* report print
* boom
* fix: numerical issues.
* fix: resolution
* comment about resize.
* change the order of the training transformation.
* save call.
* debug
* remove print
* manually detaching necessary?
* use the same vae for validation.
* add: readme.
* unwrap text encoder when saving hook only for full text encoder tuning
* unwrap text encoder when saving hook only for full text encoder tuning
* save embeddings in each checkpoint as well
* save embeddings in each checkpoint as well
* save embeddings in each checkpoint as well
* Update examples/advanced_diffusion_training/train_dreambooth_lora_sdxl_advanced.py
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* add documentation for DeepCache
* fix typo
* add wandb url for DeepCache
* fix some typos
* add item in _toctree.yml
* update formats for arguments
* Update deepcache.md
* Update docs/source/en/optimization/deepcache.md
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* add StableDiffusionXLPipeline in doc
* Separate SDPipeline and SDXLPipeline
* Add the paper link of ablation experiments for hyper-parameters
* Apply suggestions from code review
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Make WDS pipeline interpolation type configurable.
* Make the VAE encoding batch size configurable.
* Make lora_alpha and lora_dropout configurable for LCM LoRA scripts.
* Generalize scalings_for_boundary_conditions function and make the timestep scaling configurable.
* Make LoRA target modules configurable for LCM-LoRA scripts.
* Move resolve_interpolation_mode to src/diffusers/training_utils.py and make interpolation type configurable in non-WDS script.
* apply suggestions from review
* debug
* debug test_with_different_scales_fusion_equivalence
* use the right method.
* place it right.
* let's see.
* let's see again
* alright then.
* add a comment.
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
* Update unet_2d_blocks.py
similar doc-string add to have in the original diffusion repository.
* Batter way to write binarize function
* Solve check_code_quality error
* My mistake to run pull request but not reformated file
* Update image_processor.py
* remove extra variable and space
* Update image_processor.py
* Run ruff libarary to reformat my file
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* add: test to check if peft loras are loadable in non-peft envs.
* add torch_device approrpiately.
* fix: get_dummy_inputs().
* test logits.
* rename
* debug
* debug
* fix: generator
* new assertion values after fixing the seed.
* shape
* remove print statements and settle this.
* to update values.
* change values when lora config is initialized under a fixed seed.
* update colab link
* update notebook link
* sanity restored by getting the exact same values without peft.
* change timesteps used to calculate snr when --with_prior_preservation is enabled
* change timesteps used to calculate snr when --with_prior_preservation is enabled (canonical script)
* style
* revert canonical script to before snr gamma change
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Add unload_ip_adapter method
* Update attn_processors with original layers
* Add test
* Use set_default_attn_processor
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Fix gradient-checkpointing option is ignored in SDXL+LoRA training. (#6388)
* Fix gradient-checkpointing option is ignored in SD+LoRA training.
* Fix gradient checkpoint is not applied to text encoders. (SDXL+LoRA)
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* add doc for diffusion fast
* add entry to _toctree
* Apply suggestions from code review
* fix titlew
* fix: title entry
* add note about fuse_qkv_projections
* add adapter_name in fuse
* add tesrt
* up
* fix CI
* adapt from suggestion
* Update src/diffusers/utils/testing_utils.py
Co-authored-by: Benjamin Bossan <BenjaminBossan@users.noreply.github.com>
* change to `require_peft_version_greater`
* change variable names in test
* Update src/diffusers/loaders/lora.py
Co-authored-by: Benjamin Bossan <BenjaminBossan@users.noreply.github.com>
* break into 2 lines
* final comments
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Benjamin Bossan <BenjaminBossan@users.noreply.github.com>
* [Peft] fix saving / loading when unet is not "unet"
* Update src/diffusers/loaders/lora.py
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* undo stablediffusion-xl changes
* use unet_name to get unet for lora helpers
* use unet_name
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* remove validation args from textual onverson tests
* reduce number of train steps in textual inversion tests
* fix: directories.
* debig
* fix: directories.
* remove validation tests from textual onversion
* try reducing the time of test_text_to_image_checkpointing_use_ema
* fix: directories
* speed up test_text_to_image_checkpointing
* speed up test_text_to_image_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints
* fix
* speed up test_instruct_pix2pix_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints
* set checkpoints_total_limit to 2.
* test_text_to_image_lora_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints speed up
* speed up test_unconditional_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints
* debug
* fix: directories.
* speed up test_instruct_pix2pix_checkpointing_checkpoints_total_limit
* speed up: test_controlnet_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints
* speed up test_controlnet_sdxl
* speed up dreambooth tests
* speed up test_dreambooth_lora_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints
* speed up test_custom_diffusion_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints
* speed up test_text_to_image_lora_sdxl_text_encoder_checkpointing_checkpoints_total_limit
* speed up # checkpoint-2 should have been deleted
* speed up examples/text_to_image/test_text_to_image.py::TextToImage::test_text_to_image_checkpointing_checkpoints_total_limit
* additional speed ups
* style
* fix RuntimeError: Input type (float) and bias type (c10::Half) should be the same
* format source code
* format code
* remove the autocast blocks within the pipeline
* add autocast blocks to pipeline caller in train_text_to_image_lora.py
* [Community Pipeline] Add Marigold Monocular Depth Estimation
- add single-file pipeline
- update README
* fix format - add one blank line
* format script with ruff
* use direct image link in example code
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* separate out upsamplers and downsamplers.
* import all the necessary blocks in resnet for backward comp.
* move upsample2d and downsample2d to utils.
* move downsample_2d to downsamplers.py
* apply feedback
* fix import
* samplers -> sampling
* EulerAncestral add `rescale_betas_zero_snr`
Uses same infinite sigma fix from EulerDiscrete. Interestingly the
ancestral version had the opposite problem: too much contrast instead of
too little.
* UT for EulerAncestral `rescale_betas_zero_snr`
* EulerAncestral upcast samples during step()
It helps this scheduler too, particularly when the model is using bf16.
While the noise dtype is still the model's it's automatically upcasted
for the add so all it affects is determinism.
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* fix: unscale fp16 gradient problem
* fix for dreambooth lora sdxl
* make the type-casting conditional.
* Apply suggestions from code review
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix: init for vae during pixart tests
* print the values
* add flatten
* correct assertion value for test_inference
* correct assertion values for test_inference_non_square_images
* run styling
* debug test_inference_with_multiple_images_per_prompt
* fix assertion values for test_inference_with_multiple_images_per_prompt
Typo: The script for LoRA training is `train_text_to_image_lora_prior.py` not `train_text_to_image_prior_lora.py`.
Alternatively you could rename the file and keep the README.md unchanged.
* feat: introduce autoencoders module
* more changes for styling and copy fixing
* path changes in the docs.
* fix: import structure in init.
* fix controlnetxs import
* Clean up comments in LCM(-LoRA) distillation scripts.
* Calculate predicted source noise noise_pred correctly for all prediction_types.
* make style
* apply suggestions from review
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* load pipeline for inference only if validation prompt is used
* move things outside
* load pipeline for inference only if validation prompt is used
* fix readme when validation prompt is used
---------
Co-authored-by: linoytsaban <linoy@huggingface.co>
Co-authored-by: apolinário <joaopaulo.passos@gmail.com>
* fix broken example in pipeline_stable_diffusion_safe
* fix typo in pipeline_stable_diffusion_pix2pix_zero
* add missing docs
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_upscale.py
* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_attend_and_excite.py
* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_instruct_pix2pix.py
* update tests
* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_panorama.py
* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_sag.py
* support ip-adapter in src/diffusers/pipelines/stable_diffusion_safe/pipeline_stable_diffusion_safe.py
* support ip-adapter in src/diffusers/pipelines/latent_consistency_models/pipeline_latent_consistency_text2img.py
* support ip-adapter in src/diffusers/pipelines/latent_consistency_models/pipeline_latent_consistency_img2img.py
* support ip-adapter in src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_ldm3d.py
* revert changes to sd_attend_and_excite and sd_upscale
* make style
* fix broken tests
* update ip-adapter implementation to latest
* apply suggestions from review
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Fix SD scripts - there are only 2 items per batch
* Adjustments to make the SDXL scripts work with other datasets
* Use public webdataset dataset for examples
* make style
* Minor tweaks to the readmes.
* Stress that the database is illustrative.
* utils and test modifications to enable device agnostic testing
* device for manual seed in unet1d
* fix generator condition in vae test
* consistency changes to testing
* make style
* add device agnostic testing changes to source and one model test
* make dtype check fns private, log cuda fp16 case
* remove dtype checks from import utils, move to testing_utils
* adding tests for most model classes and one pipeline
* fix vae import
* Update train_dreambooth_lora_sdxl_advanced.py
* remove global function args from dreamboothdataset class
* style
* style
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* improve help tags
* style fix
* changes token_abstraction type to string.
support multiple concepts for pivotal using a comma separated string.
* style fixup
* changed logger to warning (not yet available)
* moved the token_abstraction parsing to be in the same block as where we create the mapping of identifier to token
---------
Co-authored-by: Linoy <linoy@huggingface.co>
* Update value_guided_sampling.py
Changed the scheduler step function as predict_epsilon parameter is not there in latest DDPM Scheduler
* Update value_guided_sampling.md
Updated a link to a working notebook
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* fix: duplicate unet prefix problem.
* Update src/diffusers/loaders/lora.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* adapt PixArtAlphaPipeline for pixart-lcm model
* remove original_inference_steps from __call__
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* LLMGroundedDiffusionPipeline: inherit from DiffusionPipeline and fix peft
* Use main in the revision in the examples
* Add "Copied from" statements in comments
* Fix formatting with ruff
* imports and readme bug fixes
* bug fix - ensures text_encoder params are dtype==float32 (when using pivotal tuning) even if the rest of the model is loaded in fp16
* added pivotal tuning to readme
* mapping token identifier to new inserted token in validation prompt (if used)
* correct default value of --train_text_encoder_frac
* change default value of --adam_weight_decay_text_encoder
* validation prompt generations when using pivotal tuning bug fix
* style fix
* textual inversion embeddings name change
* style fix
* bug fix - stopping text encoder optimization halfway
* readme - will include token abstraction and new inserted tokens when using pivotal tuning
- added type to --num_new_tokens_per_abstraction
* style fix
---------
Co-authored-by: Linoy Tsaban <linoy@huggingface.co>
* make `requires_safety_checker` a kwarg instead of a positional argument as it's more future-proof
* apply `make style` formatting edits
* add image_encoder to arguments and pass to super constructor
* add diffusers example
* add diffusers example
* Comment about making it faster
* Apply suggestions from code review
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
---------
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* Fixed custom module importing on Windows
Windows use back slash and `os.path.join()` follows that convention.
* Apply suggestions from code review
Co-authored-by: Lucain <lucainp@gmail.com>
* Update pipeline_utils.py
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Lucain <lucainp@gmail.com>
* integrated sdxl for the text2video-zero pipeline
* make fix-copies
* fixed CI issues
* make fix-copies
* added docs and `copied from` statements
* added fast tests
* made a small change in docs
* quality+style check fix
* updated docs. added controlnet inference with sdxl
* added device compatibility for fast tests
* fixed docstrings
* changing vae upcasting
* remove torch.empty_cache to speed up inference
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* made fast tests to run on dummy models only, fixed copied from statements
* fixed testing utils imports
* Added bullet points for SDXL support
* fixed formatting & quality
* Update tests/pipelines/text_to_video/test_text_to_video_zero_sdxl.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Update tests/pipelines/text_to_video/test_text_to_video_zero_sdxl.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fixed minor error for merging
* fixed updates of sdxl
* made fast tests inherit from `PipelineTesterMixin` and run in 3-4secs on CPU
* make style && make quality
* reimplemented fast tests w/o default attn processor
* make style & make quality
* make fix-copies
* make fix-copies
* fixed docs
* make style & make quality & make fix-copies
* bug fix in cross attention
* make style && make quality
* make fix-copies
* fix gpu issues
* make fix-copies
* updated pipeline signature
---------
Co-authored-by: Vahram <vahram.tadevosyan@lambda-loginnode02.cm.cluster>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
* Add SSD-1B support for controlnet model
* Add conditioning_channels into ControlNet init from unet
* Fix black formatting
* Isort fixes
* Adds SSD-1B controlnet pipeline test with UNetMidBlock2D as mid block
* Overrides failing ssd-1b tests
* Fixes tests after main branch update
* Fixes code quality checks
---------
Co-authored-by: Marko Kostiv <marko@linearity.io>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Change pipeline_controlnet_inpaint.py to add ip-adapter support. Changes are similar to those in pipeline_controlnet
* Change tests for the StableDiffusionControlNetInpaintPipeline by adding image_encoder: None
* Update src/diffusers/pipelines/controlnet/pipeline_controlnet_inpaint.py
Co-authored-by: YiYi Xu <yixu310@gmail.com>
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* move several state dict conversion utils out of lora.py
* check
* check
* check
* check
* check
* check
* check
* revert back
* check
* check
* again check
* maybe fix?
* Apply suggestions from code review
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* bug in MultiAdapter for Inpainting
* adapter_input is a list for MultiAdapter
---------
Co-authored-by: andres <andres@hax.ai>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* [Tests] Make sure that we don't run tests mulitple times
* [Tests] Make sure that we don't run tests mulitple times
* [Tests] Make sure that we don't run tests mulitple times
* add comments to explain the code better
* add comments to explain the code better
* add comments to explain the code better
* add comments to explain the code better
* add comments to explain the code better
* fix more
* fix more
* fix more
* fix more
* fix more
* fix more
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
* Update unet_2d_blocks.py
similar doc-string add to have in the original diffusion repository.
* I enhanced the code by replacing multiple redundant variables with a single variable, as they all served the same purpose. Additionally, I utilized the get_activation function for improved flexibility in choosing activation functions.
* Using as black package to reformated my file
* reverte some changes
* Remove conv_out_padding variables and using as conv_in_padding
* conv_out_padding create and add them into the code.
* run black command to solving styling problem
* add little bit space between comment and import statement
* I am utilizing the ruff library to address the style issues in my Makefile.
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Add custom timesteps support to LCMScheduler.
* Add custom timesteps support to StableDiffusionPipeline.
* Add custom timesteps support to StableDiffusionXLPipeline.
* Add custom timesteps support to remaining Stable Diffusion pipelines which support LCMScheduler (img2img, inpaint).
* Add custom timesteps support to remaining Stable Diffusion XL pipelines which support LCMScheduler (img2img, inpaint).
* Add custom timesteps support to StableDiffusionControlNetPipeline.
* Add custom timesteps support to T21 Stable Diffusion (XL) Adapters.
* Clean up Stable Diffusion inpaint tests.
* Manually add support for custom timesteps to AltDiffusion pipelines since make fix-copies doesn't appear to work correctly (it deletes the whole pipeline).
* make style
* Refactor pipeline timestep handling into the retrieve_timesteps function.
* deprecated: KarrasVeScheduler, ScoreSdeVpScheduler
* delete tests relevant to deprecated schedulers
* chore: run make style
* fix: import error caused due to incorrect _import_structure after deprecation
* fix: ScoreSdeVpScheduler was not importable from diffusers
* remove import added by assumption
* Update src/diffusers/schedulers/__init__.py as suggested by @patrickvonplaten
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* make it a part deprecated
* Apply suggestions from code review
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Fix
* fix
* fix doc
* fix doc....again.......
* remove karras_ve test folder
Co-Authored-By: YiYi Xu <yixu310@gmail.com>
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: yiyixuxu <yixu310@gmail,com>
* [Fix: pixart-alpha]
add ASPECT_RATIO_512_BIN in use_resolution_binning for random 512px image generation.
* add slow test file for 512px generation without resolution binning
* fix: slow tests for resolution binning.
---------
Co-authored-by: jschen <chenjunsong4@h-partners.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* finalize
* finalize
* finalize
* add slow test
* add slow test
* add slow test
* Fix more
* add slow test
* fix more
* fix more
* fix more
* fix more
* fix more
* fix more
* fix more
* fix more
* fix more
* Better
* Fix more
* Fix more
* add slow test
* Add auto pipelines
* add slow test
* Add all
* add slow test
* add slow test
* add slow test
* add slow test
* add slow test
* Apply suggestions from code review
* add slow test
* add slow test
* Additions:
- support for different lr for text encoder
- support for Prodigy optimizer
- support for min snr gamma
- support for custom captions and dataset loading from the hub
* adjusted --caption_column behaviour (to -not- use the second column of the dataset by default if --caption_column is not provided)
* fixed --output_dir / --model_dir_name confusion
* added --repeats, --adam_weight_decay_text_encoder
+ some fixes
* Update examples/dreambooth/train_dreambooth_lora_sdxl.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Update examples/dreambooth/train_dreambooth_lora_sdxl.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Update examples/dreambooth/train_dreambooth_lora_sdxl.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* - import compute_snr from diffusers/training_utils.py
- cluster adamw together
- when using 'prodigy', if --train_text_encoder == True and --text_encoder_lr != --learning rate, changes the lr of the text encoders optimization params to be --learning_rate (otherwise errors)
* shape fixes when custom captions are used
* formatting and a little cleanup
* code styling
* --repeats default value fixed, changed to 1
* bug fix - removed redundant lines of embedding concatenation when using prior_preservation (that duplicated class_prompt embeddings)
* changed dataset loading logic according to the following usecases (to avoid unnecessary dependency on datasets)-
1. user provides --dataset_name
2. user provides local dir --instance_data_dir that contains a metadata .jsonl file
3. user provides local dir --instance_data_dir that contains only images
in cases [1,2] we import datasets and use load_dataset method, in case [3] we process the data same as in the original script setting
* styling fix
* arg name fix
* adjusted the --repeats logic
* -removed redundant arg and 'if' when loading local folder with prompts
-updated readme template
-some default val fixes
-custom caption tests
* image path fix for readme
* code style
* bug fix
* --caption_column arg
* readme fix
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Linoy Tsaban <linoy@huggingface.co>
* Change LCMScheduler.set_timesteps to pick more evenly spaced inference timesteps.
* Change inference_indices implementation to better match previous behavior.
* Add num_inference_steps=26 test case to test_inference_steps.
* run CI
---------
Co-authored-by: patil-suraj <surajp815@gmail.com>
* fix an issue that ipex occupy too much memory, it will not impact performance
* make style
---------
Co-authored-by: root <jun.chen@intel.com>
Co-authored-by: Meng Guoqing <guoqing.meng@intel.com>
An upcoming change to JAX will include non-local (addressable) CPU devices in jax.devices() when JAX is used multicontroller-style, where there are multiple Python processes.
This change preserves the current behavior by replacing uses of jax.devices("cpu"), which previously only returned local devices, with jax.local_devices("cpu"), which will return local devices both now and in the future.
This change is always safe (i.e., it should always preserve the previous behavior), but it may sometimes be unnecessary if code is never used in a multicontroller setting.
Co-authored-by: Peter Hawkins <phawkins@google.com>
* fix: UnboundLocalError with image_latents
* chore: run make style, quality, fix-copies
* revert changes from make fix-copies
* revert changes from make fix-copies
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* add also peft latest on peft CI
* up
* up
* up
* Update .github/workflows/pr_test_peft_backend.yml
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* begin doc
* fix examples
* add in toctree
* fix toctree
* improve copy
* improve introductions
* add lcm doc
* fix filename
* Apply suggestions from code review
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* address Sayak's comments
* remove controlnet aux
* open in colab
* move to Specific pipeline examples
* update controlent and adapter examples
---------
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* improvement: docs and type hints
* improvement: docs and type hints
minor refactor
* improvement: docs and type hints
* update with suggestions from review
Co-Authored-By: Dhruv Nair <dhruv.nair@gmail.com>
---------
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
* Fix typos, update, add Copyright info, and trim trailing whitespace
* Update docs/source/en/api/pipelines/text_to_video_zero.md
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* 1 second is not a long video, but 6 seconds is
* Update text_to_video_zero.md
* Update text_to_video_zero.md
* Update text_to_video_zero.md
* Update wuerstchen.md
---------
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* does this fix things?
* attention mask use
* attention mask order
* better masking.
* add: tesrt
* remove mask_featur
* test
* debug
* fix: tests
* deprecate mask_feature
* add deprecation test
* add slow test
* add print statements to retrieve the assertion values.
* fix for the 1024 fast tes
* fix tesy
* fix the remaining
* Apply suggestions from code review
* more debug
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Fix the pipeline name in the examples for LMD+ pipeline
* Add LMD+ colab link
* Apply code formatting
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update the reference for text_to_video.md
The original reference (VideoFusion) might be misleading. VideoFusion is not open-sourced. I am the co-first author of ModelScopeT2V. I change the referred paper to the right one.
* Update docs/source/en/api/pipelines/text_to_video.md
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* [Docs] Running the pipeline twice does not appear to be the intention of these examples
One is with `cross_attention_kwargs` and the other (next line) removes it
* [Docs] Clarify that these are two separate examples
One using `scale` and the other without it
* add: locm docs.
* correct path
* Apply suggestions from code review
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* Apply suggestions from code review
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* Apply suggestions from code review
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* up
* add
---------
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* consistency decoder
* rename
* Apply suggestions from code review
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/pipelines/consistency_models/pipeline_consistency_models.py
* uP
* Apply suggestions from code review
* uP
* uP
* uP
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Add adapter fusing + PEFT to the docs
* Update docs/source/en/tutorials/using_peft_for_inference.md
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update docs/source/en/tutorials/using_peft_for_inference.md
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update docs/source/en/tutorials/using_peft_for_inference.md
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update docs/source/en/tutorials/using_peft_for_inference.md
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update docs/source/en/tutorials/using_peft_for_inference.md
* Update docs/source/en/tutorials/using_peft_for_inference.md
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
* Update unet_2d_blocks.py
similar doc-string add to have in the original diffusion repository.
* I removed the dummy variable defined in both the encoder and decoder.
* Now, I run black package to reformat my file
* Remove the redundant line from the adapter.py file.
* Black package using to reformated my file
* Replacing the nn.Mish activation function with a get_activation function allows developers to more easily choose the right activation function for their task. Additionally, removing redundant variables can improve code readability and maintainability.
* I try to fix this: Fast tests for PRs / Fast PyTorch Models & Schedulers CPU tests (pull_request)
* Update src/diffusers/models/resnet.py
Co-authored-by: YiYi Xu <yixu310@gmail.com>
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Refactor LCMScheduler.step such that prev_sample == denoised at the last timestep in the schedule.
* Make timestep scaling when calculating boundary conditions configurable.
* Reparameterize timestep_scaling to be a multiplicative rather than division scaling.
* make style
* fix dtype conversion
* make style
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
* Update unet_2d_blocks.py
similar doc-string add to have in the original diffusion repository.
* I removed the dummy variable defined in both the encoder and decoder.
* Now, I run black package to reformat my file
* Remove the redundant line from the adapter.py file.
* Black package using to reformated my file
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
* Update unet_2d_blocks.py
similar doc-string add to have in the original diffusion repository.
* I removed the dummy variable defined in both the encoder and decoder.
* Now, I run black package to reformat my file
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
* Update final model offload for more pipelines
Add test to ensure all pipeline components are returned to CPU after
execution with model offloading
* Add comment to explain early UNet offload in Text-to-Video pipeline
* Style
* stabilize dpmpp for sdxl by using euler at the final step
* add lu's uniform logsnr time steps
* add test
* fix check_copies
* fix tests
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix error reported 'find_unused_parameters' running in mutiple GPUs or NPUs
* fix code check of importing module by its alphabetic order
---------
Co-authored-by: jiaqiw <wangjiaqi50@huawei.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
* Update unet_2d_blocks.py
similar doc-string add to have in the original diffusion repository.
* I use a lower method in the activation function.
* Replace multiple if-else statements with a dictionary of activation functions, and call one if statement to retrieve the appropriate function.
* I am using black package to reforamted my file
* I defined the ACTIVATION_FUNCTIONS variable outside of the function
* activation function variable convert to lower case
* First, I resolved the conflict issue. Then, I ran the Black package to reformat my file.
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* improvement: add typehints and docs to src/diffusers/models/attention_processor.py
* improvement: add typehints and docs to src/diffusers/models/vae.py
* improvement: add missing docs in src/diffusers/models/vq_model.py
* improvement: add typehints and docs to src/diffusers/models/transformer_temporal.py
* improvement: add typehints and docs to src/diffusers/models/t5_film_transformer.py
* improvement: add type hints to src/diffusers/models/unet_1d_blocks.py
* improvement: add missing type hints to src/diffusers/models/unet_2d_blocks.py
* fix: CI error (make fix-copies required)
* fix: CI error (make fix-copies required again)
---------
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
* Add a new community pipeline
examples/community/latent_consistency_img2img.py
which can be called like this
import torch
from diffusers import DiffusionPipeline
pipe = DiffusionPipeline.from_pretrained(
"SimianLuo/LCM_Dreamshaper_v7", custom_pipeline="latent_consistency_txt2img", custom_revision="main")
# To save GPU memory, torch.float16 can be used, but it may compromise image quality.
pipe.to(torch_device="cuda", torch_dtype=torch.float32)
img2img=LatentConsistencyModelPipeline_img2img(
vae=pipe.vae,
text_encoder=pipe.text_encoder,
tokenizer=pipe.tokenizer,
unet=pipe.unet,
#scheduler=pipe.scheduler,
scheduler=None,
safety_checker=None,
feature_extractor=pipe.feature_extractor,
requires_safety_checker=False,
)
img = Image.open("thisismyimage.png")
result = img2img(prompt,img,strength,num_inference_steps=4)
* Apply suggestions from code review
Fix name formatting for scheduler
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* update readme (and run formatter on latent_consistency_img2img.py)
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix
* fix copies
* remove heun from tests
* add back heun and fix the tests to include 2nd order
* fix the other test too
* Apply suggestions from code review
* Apply suggestions from code review
* Apply suggestions from code review
* make style
* add more comments
---------
Co-authored-by: yiyixuxu <yixu310@gmail,com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* initial commit for LatentConsistencyModelPipeline and LCMScheduler based on the community pipeline
* Add callback and freeu support.
* apply suggestions from review
* Clean up LCMScheduler
* Remove timeindex argument to LCMScheduler.step.
* Add support for clipping or thresholding the predicted original sample.
* Remove unused methods and arguments in LCMScheduler.
* Improve comment about (lack of) negative prompt support.
* Change input guidance_scale to match the StableDiffusionPipeline (Imagen) CFG formulation.
* Move lcm_origin_steps from pipeline __call__ to LCMScheduler.__init__/config (as origin_steps).
* Fix typo when clipping/thresholding in LCMScheduler.
* Add some initial LCMScheduler tests.
* add type annotations from review
* Fix type annotation bug.
* Override test_add_noise_device in LCMSchedulerTest since hardcoded timesteps doesn't work under default settings.
* Add generator argument pipeline prepare_latents call.
* Cast LCMScheduler.timesteps to long in set_timesteps.
* Add onestep and multistep full loop scheduler tests.
* Set default height/width to None and don't hardcode guidance scale embedding dim.
* Add initial LatentConsistencyPipeline fast and slow tests.
* Add initial documentation for LatentConsistencyModelPipeline and LCMScheduler.
* Make remaining failing fast tests pass.
* make style
* Make original_inference_steps configurable from pipeline __call__ again.
* make style
* Remove guidance_rescale arg from pipeline __call__ since LCM currently doesn't support CFG.
* Make LCMScheduler defaults match config of LCM_Dreamshaper_v7 checkpoint.
* Fix LatentConsistencyPipeline slow tests and add dummy expected slices.
* Add checks for original_steps in LCMScheduler.set_timesteps.
* make fix-copies
* Improve LatentConsistencyModelPipeline docs.
* Apply suggestions from code review
Co-authored-by: Aryan V S <avs050602@gmail.com>
* Apply suggestions from code review
Co-authored-by: Aryan V S <avs050602@gmail.com>
* Apply suggestions from code review
Co-authored-by: Aryan V S <avs050602@gmail.com>
* Update src/diffusers/schedulers/scheduling_lcm.py
* Apply suggestions from code review
Co-authored-by: Aryan V S <avs050602@gmail.com>
* finish
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Aryan V S <avs050602@gmail.com>
* add
* Update docs/source/en/api/pipelines/controlnet_sdxl.md
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
---------
Co-authored-by: yiyixuxu <yixu310@gmail,com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Update get_dummy_inputs(...) in T2I-Adapter tests to take image height and width as params.
* Update the T2I-Adapter unit tests to run with the standard number of UNet down blocks so that all T2I-Adapter down blocks get exercised.
* Update the T2I-Adapter down blocks to better match the padding behavior of the UNet.
* Revert "Update the T2I-Adapter unit tests to run with the standard number of UNet down blocks so that all T2I-Adapter down blocks get exercised."
This reverts commit 6d4a060a34.
* Create utility functions for testing the T2I-Adapter downscaling bahevior.
* (minor) Improve readability with an intermediate named variable.
* Statically parameterize T2I-Adapter test dimensions rather than generating them dynamically.
* Fix static checks.
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Added args, kwargs to ```U
* Add UNetMidBlock2D as a supported mid block type
* Fix extra init input for UNetMidBlock2D, change allowed types for Mid-block init
* Update unet_2d_condition.py
* Update unet_2d_condition.py
* Update unet_2d_condition.py
* Update unet_2d_condition.py
* Update unet_2d_condition.py
* Update unet_2d_condition.py
* Update unet_2d_condition.py
* Update unet_2d_condition.py
* Update unet_2d_blocks.py
* Update unet_2d_blocks.py
* Update unet_2d_blocks.py
* Update unet_2d_condition.py
* Update unet_2d_blocks.py
* Updated docstring, increased check strictness
Updated the docstring for ```UNet2DConditionModel``` to include ```reverse_transformer_layers_per_block``` and updated checking for nested list type ```transformer_layers_per_block```
* Add basic shape-check test for asymmetrical unets
* Update src/diffusers/models/unet_2d_blocks.py
Removed blank line
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_condition.py
Remove blank space
* Update unet_2d_condition.py
Changed docstring for `mid_block_type`
* Fixed docstring and wrong default value
* Reformat with black
* Reformat with necessary commands
* Add UNetMidBlockFlat to versatile_diffusion/modeling_text_unet.py to ensure consistency
* Removed args, kwargs, use on mid-block type
* Make fix-copies
* Update src/diffusers/models/unet_2d_condition.py
Wrap into single line
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* make fix-copies
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
* Update unet_2d_blocks.py
similar doc-string add to have in the original diffusion repository.
* Update unet_2d_blocks.py
Added Beutifull doc-string into the UNetMidBlock2D class.
* Update unet_2d_blocks.py
I replaced the definition in this parameter resnet_time_scale_shift and resnet_groups.
* Update unet_2d_blocks.py
I remove additional sentences into the resnet_groups argument.
* Update unet_2d_blocks.py
I replaced my definition with the maintainer definition in the attention_head_dim parameter.
* I am using black package for reformated my file
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* added TODOs
* Enhanced and reformatted the docstrings of IFPipeline methods.
* Enhanced and fixed the docstrings of IFImg2ImgSuperResolutionPipeline methods.
* Enhanced and fixed the docstrings of IFImg2ImgPipeline methods.
* Enhanced and fixed the docstrings of IFInpaintingSuperResolutionPipeline methods.
* Enhanced and fixed the docstrings of IFInpaintingPipeline methods.
* Enhanced and fixed the docstrings of IFSuperResolutionPipeline methods.
* Update src/diffusers/pipelines/deepfloyd_if/pipeline_if.py
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update src/diffusers/pipelines/deepfloyd_if/pipeline_if_img2img.py
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update src/diffusers/pipelines/deepfloyd_if/pipeline_if_img2img_superresolution.py
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update src/diffusers/pipelines/deepfloyd_if/pipeline_if_inpainting.py
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update src/diffusers/pipelines/deepfloyd_if/pipeline_if_superresolution.py
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update src/diffusers/pipelines/deepfloyd_if/pipeline_if_inpainting_superresolution.py
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* remove redundant code
* fix code style
* revert the ordering to not break backwards compatibility
---------
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* changed channel parameters for UNET and VAE. Decreased hidden layers size with increased attention heads and intermediate size
* changed the assertion check range
* clean up
---------
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
* fix: sdxl pipeline when unet is not available.
* fix moe
* account for text
* ifx more
* don't make unet optional.
* Apply suggestions from code review
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* split conditionals.
* add optional components to sdxl pipeline
* propagate changes to the rest of the pipelines.
* add: test
* add to all
* fix: rest of the pipelines.
* use pipeline_class variable
* separate pipeline mixin
* use safe_serialization
* fix: test
* access actual output.
* add: optional test to adapter and ip2p sdxl pipeline tests/
* add optional test to controlnet sdxl.
* fix tests
* fix ip2p tests
* fix more
* fifx more.
* use np output type.
* fix for StableDiffusionXLMultiControlNetPipelineFastTests.
* fix: SDXLOptionalComponentsTesterMixin
* Apply suggestions from code review
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix tests
* Empty-Commit
* revert previous
* quality
* fix: test
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Add ability to mix usage of T2I-Adapter(s) and ControlNet(s).
Previously, UNet2DConditional implemnetation onloy allowed use of one or the other.
Adds new forward() arg down_intrablock_additional_residuals specifically for T2I-Adapters. If down_intrablock_addtional_residuals is not used, maintains backward compatibility with prior usage of only T2I-Adapter or ControlNet but not both
* Improving forward() arg docs in src/diffusers/models/unet_2d_condition.py
Co-authored-by: psychedelicious <4822129+psychedelicious@users.noreply.github.com>
* Add deprecation warning if down_block_additional_residues is used for T2I-Adapter (intrablock residuals)
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Oops my bad, fixing last commit.
* Added import of diffusers utils.deprecate
* Conform to max line length
* Modifying T2I-Adapter pipelines to reflect change to UNet forward() arg for T2I-Adapter residuals.
---------
Co-authored-by: psychedelicious <4822129+psychedelicious@users.noreply.github.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* add: freeu to the core sdxl pipeline.
* add: freeu to video2video
* add: freeu to the core SD pipelines.
* add: freeu to image variation for sdxl.
* add: freeu to SD ControlNet pipelines.
* add: freeu to SDXL controlnet pipelines.
* add: freu to t2i adapter pipelines.
* make fix-copies.
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* improvement: add missing typehints and docs to diffusers/models/attention.py
* chore: convert doc strings to raw python strings
add missing typehints
* improvement: add missing typehints and docs to diffusers/models/adapter.py
* improvement: add missing typehints and docs to diffusers/models/lora.py
* docs: include suggestion by @sayakpaul in src/diffusers/models/adapter.py
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* docs: include suggestion by @sayakpaul in src/diffusers/models/adapter.py
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* docs: include suggestion by @sayakpaul in src/diffusers/models/adapter.py
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* docs: include suggestion by @sayakpaul in src/diffusers/models/adapter.py
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/lora.py
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Added mark_step for sdxl to run with pytorch xla. Also updated README with instructions for xla
* adding soft dependency on torch_xla
* fix some styling
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
* add missing docstrings
* chore: run make quality
* improvement: include docs suggestion by @yiyixuxu
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* decrease UNet2DConditionModel & ControlNetModel blocks
* decrease UNet2DConditionModel & ControlNetModel blocks
* decrease even more blocks & number of norm groups
* decrease vae block out channels and n of norm goups
* fix code style
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* fix(gligen_inpaint_pipeline): 🐛 Wrap the timestep() 0-d tensor in a list to convert to 1-d tensor. This avoids the TypeError caused by trying to directly iterate over a 0-dimensional tensor in the denoising stage
* test(gligen/gligen_text_image): unit test using the EulerAncestralDiscreteScheduler
---------
Co-authored-by: zhen-hao.chu <zhen-hao.chu@vitrox.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Min-SNR Gamma: correct the fix for SNR weighted loss in v-prediction by adding 1 to SNR rather than the resulting loss weights
Co-authored-by: bghira <bghira@users.github.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* ✨ Added Fourier filter function to upsample blocks
* 🔧 Update Fourier_filter for float16 support
* ✨ Added UNetFreeUConfig to UNet model for FreeU adaptation 🛠️
* move unet to its original form and add fourier_filter to torch_utils.
* implement freeU enable mechanism
* implement disable mechanism
* resolution index.
* correct resolution idx condition.
* fix copies.
* no need to use resolution_idx in vae.
* spell out the kwargs
* proper config property
* fix attribution setting
* place unet hasattr properly.
* fix: attribute access.
* proper disable
* remove validation method.
* debug
* debug
* debug
* debug
* debug
* debug
* potential fix.
* add: doc.
* fix copies
* add: tests.
* add: support freeU in SDXL.
* set default value of resolution idx.
* set default values for resolution_idx.
* fix copies
* fix rest.
* fix copies
* address PR comments.
* run fix-copies
* move apply_free_u to utils and other minors.
* introduce support for video (unet3D)
* minor ups
* consistent fix-copies.
* consistent stuff
* fix-copies
* add: rest
* add: docs.
* fix: tests
* fix: doc path
* Apply suggestions from code review
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* style up
* move to techniques.
* add: slow test for sd freeu.
* add: slow test for sd freeu.
* add: slow test for sd freeu.
* add: slow test for sd freeu.
* add: slow test for sd freeu.
* add: slow test for sd freeu.
* add: slow test for video with freeu
* add: slow test for video with freeu
* add: slow test for video with freeu
* style
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* handle case when controlnet is list
* Update src/diffusers/loaders.py
* Apply suggestions from code review
* Update src/diffusers/loaders.py
* typecheck comment
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* pipline fetcher
* update script
* clean up
* clean up
* clean up
* new pipeline runner
* rename tests to match modules
* test actions in pr
* change runner to gpu
* clean up
* clean up
* clean up
* fix report
* fix reporting
* clean up
* show test stats in failure reports
* give names to jobs
* add lora tests
* split torch cuda tests and add compile tests
* clean up
* fix tests
* change push to run only on main
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Update Unipc einsum to support 1D and 3D diffusion.
* Add unittest
* Update unittest & edge case
* Fix unittest
* Fix testing_utils.py
* Fix unittest file
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Add docstring for the AutoencoderKL's encode
#5229
* Support Python 3.8 syntax in AutoencoderKL.decode type hints
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Follow the style guidelines in AutoencoderKL's encode
#5230
---------
Co-authored-by: stano <>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Add VAE slicing and tiling methods.
* Switch to using VaeImageProcessing for preprocessing and postprocessing of images.
* Rename the VaeImageProcessor to vae_image_processor to avoid a name clash with the CLIPImageProcessor (image_processor).
* Remove the postprocess() function because we're using a VaeImageProcessor instead.
* Remove UniDiffuserPipeline.decode_image_latents because we're using VaeImageProcessor instead.
* Refactor generating text from text latents into a decode_text_latents method.
* Add enable_full_determinism() to UniDiffuser tests.
* make style
* Add PipelineLatentTesterMixin to UniDiffuserPipelineFastTests.
* Remove enable_model_cpu_offload since it is now part of DiffusionPipeline.
* Rename the VaeImageProcessor instance to self.image_processor for consistency with other pipelines and rename the CLIPImageProcessor instance to clip_image_processor to avoid a name clash.
* Update UniDiffuser conversion script.
* Make safe_serialization configurable in UniDiffuser conversion script.
* Rename image_processor to clip_image_processor in UniDiffuser tests.
* Add PipelineKarrasSchedulerTesterMixin to UniDiffuserPipelineFastTests.
* Add initial test for compiling the UniDiffuser model (not tested yet).
* Update encode_prompt and _encode_prompt to match that of StableDiffusionPipeline.
* Turn off standard classifier-free guidance for now.
* make style
* make fix-copies
* apply suggestions from review
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* added docstrings in forward methods of T2IAdapter model and FullAdapter model
* added docstrings in forward methods of FullAdapterXL and AdapterBlock models
* Added docstrings in forward methods of adapter models
* fix ddim inverse scheduler
* update test of ddim inverse scheduler
* update test of pix2pix_zero
* update test of diffedit
* fix typo
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* split_head_dim flax attn
* Make split_head_dim non default
* make style and make quality
* add description for split_head_dim flag
* Update src/diffusers/models/attention_flax.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
---------
Co-authored-by: Juan Acevedo <jfacevedo@google.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Timestep bias for fine-tuning SDXL
* Adjust parameter choices to include "range" and reword the help statements
* Condition our use of weighted timesteps on the value of timestep_bias_strategy
* style
---------
Co-authored-by: bghira <bghira@users.github.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Fix FullAdapterXL.total_downscale_factor.
* Fix incorrect error message in T2IAdapter.__init__(...).
* Move IP-Adapter test_total_downscale_factor(...) to pipeline test file (requested in code review).
* Add more info to error message about an unsupported T2I-Adapter adapter_type.
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Make sure the repo_id is valid before sending it to huggingface_hub to get a more understandable error message.
Re #5110
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* SDXL microconditioning documentation should indicate the correct default order of parameters, so that developers know
* SDXL microconditioning documentation should indicate the correct default order of parameters, so that developers know
* empty
---------
Co-authored-by: bghira <bghira@users.github.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* support transformer_layers_per block in flax UNet
* add support for text_time additional embeddings to Flax UNet
* rename attention layers for VAE
* add shape asserts when renaming attention layers
* transpose VAE attention layers
* add pipeline flax SDXL code [WIP]
* continue add pipeline flax SDXL code [WIP]
* cleanup
* Working on JIT support
Fixed prompt embedding shapes so they work in parallel mode. Assuming we
always have both text encoders for now, for simplicity.
* Fixing embeddings (untested)
* Remove spurious line
* Shard guidance_scale when jitting.
* Decode images
* Fix sharding
* style
* Refiner UNet can be loaded.
* Refiner / img2img pipeline
* Allow latent outputs from base and latent inputs in refiner
This makes it possible to chain base + refiner without having to use the
vae decoder in the base model, the vae encoder in the refiner, skipping
conversions to/from PIL, and avoiding TPU <-> CPU memory copies.
* Adapt to FlaxCLIPTextModelOutput
* Update Flax XL pipeline to FlaxCLIPTextModelOutput
* make fix-copies
* make style
* add euler scheduler
* Fix import
* Fix copies, comment unused code.
* Fix SDXL Flax imports
* Fix euler discrete begin
* improve init import
* finish
* put discrete euler in init
* fix flax euler
* Fix more
* make style
* correct init
* correct init
* Temporarily remove FlaxStableDiffusionXLImg2ImgPipeline
* correct pipelines
* finish
---------
Co-authored-by: Martin Müller <martin.muller.me@gmail.com>
Co-authored-by: patil-suraj <surajp815@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* min-SNR gamma for Dreambooth training
* Align the mse_loss_weights style with SDXL training example
---------
Co-authored-by: bghira <bghira@users.github.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Resolve v_prediction issue for min-SNR gamma weighted loss function
* Combine MSE loss calculation of epsilon and velocity, with a note about the application of the epsilon code to sample prediction
* style
---------
Co-authored-by: bghira <bghira@users.github.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* fix test
* initial commit
* change test
* updates:
* fix tests
* test fix
* test fix
* fix tests
* make test faster
* clean up
* fix precision in test
* fix precision
* Fix tests
* Fix logging test
* fix test
* fix test
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* [SDXL] Make sure multi batch prompt embeds works
* [SDXL] Make sure multi batch prompt embeds works
* improve more
* improve more
* Apply suggestions from code review
Fixed `get_word_inds` mistake/typo in P2P community pipeline
The function `get_word_inds` was taking a string of text and either a word (str) or a word index (int) and returned the indices of token(s) the word would be encoded to.
However, there was a typo, in which in the second `if` branch the word was checked to be a `str` **again**, not `int`, which resulted in an [example code from the docs](https://github.com/huggingface/diffusers/tree/main/examples/community#prompt2prompt-pipeline) to result in an error
* add support for clip skip
* fix condition
* fix
* add clip_output_layer_to_default
* expose
* remove the previous functions.
* correct condition.
* apply final layer norm
* address feedback
* Apply suggestions from code review
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* refactor clip_skip.
* port to the other pipelines.
* fix copies one more time
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Remove logger.info statement from Unet2DCondition code to ensure torch compile reliably succeeds
* Convert logging statement to a comment for future archaeologists
* Update src/diffusers/models/unet_2d_condition.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
---------
Co-authored-by: bghira <bghira@users.github.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Add attn_groups argument to UNet2DMidBlock2D to control theinternal Attention block's GroupNorm.
* Add docstring for attn_norm_num_groups in UNet2DModel.
* Since the test UNet config uses resnet_time_scale_shift == 'scale_shift', also set attn_norm_num_groups to 32.
* Add test for attn_norm_num_groups to UNet2DModelTests.
* Fix expected slices for slow tests.
* Also fix tolerances for slow tests.
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Initial commit P2P
* Replaced CrossAttention, added test skeleton
* bug fixes
* Updated docstring
* Removed unused function
* Created tests
* improved tests
- made fast inference tests faster
- corrected image shape assertions
* Corrected expected output shape in tests
* small fix: test inputs
* Update tests
- used conditional unet2d
- set expected image slices
- edit_kwargs are now not popped, so pipe can be run multiple times
* Fixed bug in int tests
* Fixed tests
* Linting
* Create prompt2prompt.md
* Added to docs toc
* Ran make fix-copies
* Fixed code blocks in docs
* Using same interface as StableDiffusionPipeline
* Fixed small test bug
* Added all options SDPipeline.__call_ has
* Fixed docstring; made __call__ like in SD
* Linting
* Added test for multiple prompts
* Improved docs
* Incorporated feedback
* Reverted formatting on unrelated files
* Moved prompt2prompt to community
- Moved prompt2prompt pipeline from main to community
- Deleted tests
- Moved documentation to community and shorted it
* Update src/diffusers/utils/dummy_torch_and_transformers_objects.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* check out dtypes.
* check out dtypes.
* check out dtypes.
* check out dtypes.
* check out dtypes.
* check out dtypes.
* check out dtypes.
* potential fix
* check out dtypes.
* check out dtypes.
* working?
Thus, issues are of the same importance as pull requests when contributing to this library ❤️.
In order to make your issue as **useful for the community as possible**, let's try to stick to some simple guidelines:
- 1. Please try to be as precise and concise as possible.
*Give your issue a fitting title. Assume that someone which very limited knowledge of diffusers can understand your issue. Add links to the source code, documentation other issues, pull requests etc...*
*Give your issue a fitting title. Assume that someone which very limited knowledge of Diffusers can understand your issue. Add links to the source code, documentation other issues, pull requests etc...*
- 2. If your issue is about something not working, **always** provide a reproducible code snippet. The reader should be able to reproduce your issue by **only copy-pasting your code snippet into a Python shell**.
*The community cannot solve your issue if it cannot reproduce it. If your bug is related to training, add your training script and make everything needed to train public. Otherwise, just add a simple Python code snippet.*
- 3. Add the **minimum amount of code / context that is needed to understand, reproduce your issue**.
- 3. Add the **minimum** amount of code / context that is needed to understand, reproduce your issue.
*Make the life of maintainers easy. `diffusers` is getting many issues every day. Make sure your issue is about one bug and one bug only. Make sure you add only the context, code needed to understand your issues - nothing more. Generally, every issue is a way of documenting this library, try to make it a good documentation entry.*
- 4. For issues related to community pipelines (i.e., the pipelines located in the `examples/community` folder), please tag the author of the pipeline in your issue thread as those pipelines are not maintained.
- type:markdown
attributes:
value:|
For more in-detail information on how to write good issues you can have a look [here](https://huggingface.co/course/chapter8/5?fw=pt)
For more in-detail information on how to write good issues you can have a look [here](https://huggingface.co/course/chapter8/5?fw=pt).
- type:textarea
id:bug-description
attributes:
@@ -46,7 +47,7 @@ body:
attributes:
label:System Info
description:Please share your system info with us. You can run the command `diffusers-cli env` and copy-paste its output below.
about: Start a new translation effort in your language
title: '[<languageCode>] Translating docs to <languageName>'
labels: WIP
assignees: ''
---
<!--
Note: Please search to see if an issue already exists for the language you are trying to translate.
-->
Hi!
Let's bring the documentation to all the <languageName>-speaking community 🌐.
Who would want to translate? Please follow the 🤗 [TRANSLATING guide](https://github.com/huggingface/diffusers/blob/main/docs/TRANSLATING.md). Here is a list of the files ready for translation. Let us know in this issue if you'd like to translate any, and we'll add your name to the list.
Some notes:
* Please translate using an informal tone (imagine you are talking with a friend about Diffusers 🤗).
* Please translate in a gender-neutral way.
* Add your translations to the folder called `<languageCode>` inside the [source folder](https://github.com/huggingface/diffusers/tree/main/docs/source).
* Register your translation in `<languageCode>/_toctree.yml`; please follow the order of the [English version](https://github.com/huggingface/diffusers/blob/main/docs/source/en/_toctree.yml).
* Once you're finished, open a pull request and tag this issue by including #issue-number in the description, where issue-number is the number of this issue. Please ping @stevhliu for review.
* 🙋 If you'd like others to help you with the translation, you can also post in the 🤗 [forums](https://discuss.huggingface.co/c/discussion-related-to-httpsgithubcomhuggingfacediffusers/63).
- [ ] This PR fixes a typo or improves the docs (you can dismiss the other checks if that's the case).
- [ ] Did you read the [contributor guideline](https://github.com/huggingface/diffusers/blob/main/CONTRIBUTING.md)?
- [ ] Did you read our [philosophy doc](https://github.com/huggingface/diffusers/blob/main/PHILOSOPHY.md) (important for complex PRs)?
- [ ] Was this discussed/approved via a Github issue or the [forum](https://discuss.huggingface.co/)? Please add a link to it if that's the case.
- [ ] Was this discussed/approved via a GitHub issue or the [forum](https://discuss.huggingface.co/c/discussion-related-to-httpsgithubcomhuggingfacediffusers/63)? Please add a link to it if that's the case.
- [ ] Did you make sure to update the documentation with your changes? Here are the
[documentation guidelines](https://github.com/huggingface/diffusers/tree/main/docs), and
[here are tips on formatting docstrings](https://github.com/huggingface/transformers/tree/main/docs#writing-source-documentation).
[here are tips on formatting docstrings](https://github.com/huggingface/diffusers/tree/main/docs#writing-source-documentation).
- [ ] Did you write any new necessary tests?
@@ -31,20 +31,20 @@ Fixes # (issue)
Anyone in the community is free to review the PR once the tests have passed. Feel free to tag
members/contributors who may be interested in your PR.
<!-- Your PR will be replied to more quickly if you can figure out the right person to tag with @
<!-- Your PR will be replied to more quickly if you can figure out the right person to tag with @.
If you know how to use git blame, that is the easiest way, otherwise, here is a rough guide of **who to tag**.
Please tag fewer than 3 people.
Core library:
- Schedulers: @williamberman and @patrickvonplaten
- Pipelines: @patrickvonplaten and @sayakpaul
- Training examples: @sayakpauland @patrickvonplaten
- Docs: @stevhliu and @yiyixuxu
- Schedulers: @yiyixuxu
- Pipelines: @sayakpaul@yiyixuxu@DN6
- Training examples: @sayakpaul
- Docs: @stevhliu and @sayakpaul
- JAX and MPS: @pcuenca
- Audio: @sanchit-gandhi
- General functionalities: @patrickvonplaten and @sayakpaul
- General functionalities: @sayakpaul@yiyixuxu@DN6
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -14,7 +14,7 @@ specific language governing permissions and limitations under the License.
We ❤️ contributions from the open-source community! Everyone is welcome, and all types of participation –not just code– are valued and appreciated. Answering questions, helping others, reaching out, and improving the documentation are all immensely valuable to the community, so don't be afraid and get involved if you're up for it!
Everyone is encouraged to start by saying 👋 in our public Discord channel. We discuss the latest trends in diffusion models, ask questions, show off personal projects, help each other with contributions, or just hang out ☕. <a href="https://Discord.gg/G7tWnz98XR"><img alt="Join us on Discord" src="https://img.shields.io/Discord/823813159592001537?color=5865F2&logo=Discord&logoColor=white"></a>
Everyone is encouraged to start by saying 👋 in our public Discord channel. We discuss the latest trends in diffusion models, ask questions, show off personal projects, help each other with contributions, or just hang out ☕. <a href="https://discord.gg/G7tWnz98XR"><img alt="Join us on Discord" src="https://img.shields.io/discord/823813159592001537?color=5865F2&logo=Discord&logoColor=white"></a>
Whichever way you choose to contribute, we strive to be part of an open, welcoming, and kind community. Please, read our [code of conduct](https://github.com/huggingface/diffusers/blob/main/CODE_OF_CONDUCT.md) and be mindful to respect it during your interactions. We also recommend you become familiar with the [ethical guidelines](https://huggingface.co/docs/diffusers/conceptual/ethical_guidelines) that guide our project and ask you to adhere to the same principles of transparency and responsibility.
@@ -28,11 +28,11 @@ the core library.
In the following, we give an overview of different ways to contribute, ranked by difficulty in ascending order. All of them are valuable to the community.
* 1. Asking and answering questions on [the Diffusers discussion forum](https://discuss.huggingface.co/c/discussion-related-to-httpsgithubcomhuggingfacediffusers) or on [Discord](https://discord.gg/G7tWnz98XR).
* 2. Opening new issues on [the GitHub Issues tab](https://github.com/huggingface/diffusers/issues/new/choose)
* 3. Answering issues on [the GitHub Issues tab](https://github.com/huggingface/diffusers/issues)
* 2. Opening new issues on [the GitHub Issues tab](https://github.com/huggingface/diffusers/issues/new/choose).
* 3. Answering issues on [the GitHub Issues tab](https://github.com/huggingface/diffusers/issues).
* 4. Fix a simple issue, marked by the "Good first issue" label, see [here](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22good+first+issue%22).
* 5. Contribute to the [documentation](https://github.com/huggingface/diffusers/tree/main/docs/source).
* 6. Contribute a [Community Pipeline](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3Acommunity-examples)
* 6. Contribute a [Community Pipeline](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3Acommunity-examples).
* 7. Contribute to the [examples](https://github.com/huggingface/diffusers/tree/main/examples).
* 8. Fix a more difficult issue, marked by the "Good second issue" label, see [here](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22Good+second+issue%22).
* 9. Add a new pipeline, model, or scheduler, see ["New Pipeline/Model"](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22New+pipeline%2Fmodel%22) and ["New scheduler"](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22New+scheduler%22) issues. For this contribution, please have a look at [Design Philosophy](https://github.com/huggingface/diffusers/blob/main/PHILOSOPHY.md).
@@ -40,7 +40,7 @@ In the following, we give an overview of different ways to contribute, ranked by
As said before, **all contributions are valuable to the community**.
In the following, we will explain each contribution a bit more in detail.
For all contributions 4.-9. you will need to open a PR. It is explained in detail how to do so in [Opening a pull requst](#how-to-open-a-pr)
For all contributions 4-9, you will need to open a PR. It is explained in detail how to do so in [Opening a pull request](#how-to-open-a-pr).
### 1. Asking and answering questions on the Diffusers discussion forum or on the Diffusers Discord
@@ -63,7 +63,7 @@ In the same spirit, you are of immense help to the community by answering such q
**Please** keep in mind that the more effort you put into asking or answering a question, the higher
the quality of the publicly documented knowledge. In the same way, well-posed and well-answered questions create a high-quality knowledge database accessible to everybody, while badly posed questions or answers reduce the overall quality of the public knowledge database.
In short, a high quality question or answer is *precise*, *concise*, *relevant*, *easy-to-understand*, *accesible*, and *well-formated/well-posed*. For more information, please have a look through the [How to write a good issue](#how-to-write-a-good-issue) section.
In short, a high quality question or answer is *precise*, *concise*, *relevant*, *easy-to-understand*, *accessible*, and *well-formated/well-posed*. For more information, please have a look through the [How to write a good issue](#how-to-write-a-good-issue) section.
**NOTE about channels**:
[*The forum*](https://discuss.huggingface.co/c/discussion-related-to-httpsgithubcomhuggingfacediffusers/63) is much better indexed by search engines, such as Google. Posts are ranked by popularity rather than chronologically. Hence, it's easier to look up questions and answers that we posted some time ago.
@@ -91,12 +91,12 @@ open a new issue nevertheless and link to the related issue.
New issues usually include the following.
#### 2.1. Reproducible, minimal bug reports.
#### 2.1. Reproducible, minimal bug reports
A bug report should always have a reproducible code snippet and be as minimal and concise as possible.
This means in more detail:
- Narrow the bug down as much as you can, **do not just dump your whole code file**
- Format your code
- Narrow the bug down as much as you can, **do not just dump your whole code file**.
- Format your code.
- Do not include any external libraries except for Diffusers depending on them.
- **Always** provide all necessary information about your environment; for this, you can run: `diffusers-cli env` in your shell and copy-paste the displayed information to the issue.
- Explain the issue. If the reader doesn't know what the issue is and why it is an issue, she cannot solve it.
@@ -105,9 +105,9 @@ This means in more detail:
For more information, please have a look through the [How to write a good issue](#how-to-write-a-good-issue) section.
You can open a bug report [here](https://github.com/huggingface/diffusers/issues/new/choose).
You can open a bug report [here](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=bug&projects=&template=bug-report.yml).
#### 2.2. Feature requests.
#### 2.2. Feature requests
A world-class feature request addresses the following points:
@@ -125,21 +125,21 @@ Awesome! Tell us what problem it solved for you.
You can open a feature request [here](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feature_request.md&title=).
#### 2.3 Feedback.
#### 2.3 Feedback
Feedback about the library design and why it is good or not good helps the core maintainers immensely to build a user-friendly library. To understand the philosophy behind the current design philosophy, please have a look [here](https://huggingface.co/docs/diffusers/conceptual/philosophy). If you feel like a certain design choice does not fit with the current design philosophy, please explain why and how it should be changed. If a certain design choice follows the design philosophy too much, hence restricting use cases, explain why and how it should be changed.
If a certain design choice is very useful for you, please also leave a note as this is great feedback for future design decisions.
You can open an issue about feedback [here](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feedback.md&title=).
#### 2.4 Technical questions.
#### 2.4 Technical questions
Technical questions are mainly about why certain code of the library was written in a certain way, or what a certain part of the code does. Please make sure to link to the code in question and please provide detail on
why this part of the code is difficult to understand.
You can open an issue about a technical question [here](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=bug&template=bug-report.yml).
#### 2.5 Proposal to add a new model, scheduler, or pipeline.
#### 2.5 Proposal to add a new model, scheduler, or pipeline
If the diffusion model community released a new model, pipeline, or scheduler that you would like to see in the Diffusers library, please provide the following information:
@@ -156,19 +156,19 @@ You can open a request for a model/pipeline/scheduler [here](https://github.com/
Answering issues on GitHub might require some technical knowledge of Diffusers, but we encourage everybody to give it a try even if you are not 100% certain that your answer is correct.
Some tips to give a high-quality answer to an issue:
- Be as concise and minimal as possible
- Be as concise and minimal as possible.
- Stay on topic. An answer to the issue should concern the issue and only the issue.
- Provide links to code, papers, or other sources that prove or encourage your point.
- Answer in code. If a simple code snippet is the answer to the issue or shows how the issue can be solved, please provide a fully reproducible code snippet.
Also, many issues tend to be simply off-topic, duplicates of other issues, or irrelevant. It is of great
help to the maintainers if you can answer such issues, encouraging the author of the issue to be
more precise, provide the link to a duplicated issue or redirect them to [the forum](https://discuss.huggingface.co/c/discussion-related-to-httpsgithubcomhuggingfacediffusers/63) or [Discord](https://discord.gg/G7tWnz98XR)
more precise, provide the link to a duplicated issue or redirect them to [the forum](https://discuss.huggingface.co/c/discussion-related-to-httpsgithubcomhuggingfacediffusers/63) or [Discord](https://discord.gg/G7tWnz98XR).
If you have verified that the issued bug report is correct and requires a correction in the source code,
please have a look at the next sections.
For all of the following contributions, you will need to open a PR. It is explained in detail how to do so in the [Opening a pull requst](#how-to-open-a-pr) section.
For all of the following contributions, you will need to open a PR. It is explained in detail how to do so in the [Opening a pull request](#how-to-open-a-pr) section.
### 4. Fixing a "Good first issue"
@@ -202,7 +202,7 @@ Please have a look at [this page](https://github.com/huggingface/diffusers/tree/
### 6. Contribute a community pipeline
[Pipelines](https://huggingface.co/docs/diffusers/api/pipelines/overview) are usually the first point of contact between the Diffusers library and the user.
Pipelines are examples of how to use Diffusers [models](https://huggingface.co/docs/diffusers/api/models) and [schedulers](https://huggingface.co/docs/diffusers/api/schedulers/overview).
Pipelines are examples of how to use Diffusers [models](https://huggingface.co/docs/diffusers/api/models/overview) and [schedulers](https://huggingface.co/docs/diffusers/api/schedulers/overview).
We support two types of pipelines:
- Official Pipelines
@@ -242,27 +242,27 @@ We support two types of training examples:
Research training examples are located in [examples/research_projects](https://github.com/huggingface/diffusers/tree/main/examples/research_projects) whereas official training examples include all folders under [examples](https://github.com/huggingface/diffusers/tree/main/examples) except the `research_projects` and `community` folders.
The official training examples are maintained by the Diffusers' core maintainers whereas the research training examples are maintained by the community.
This is because of the same reasons put forward in [6. Contribute a community pipeline](#contribute-a-community-pipeline) for official pipelines vs. community pipelines: It is not feasible for the core maintainers to maintain all possible training methods for diffusion models.
This is because of the same reasons put forward in [6. Contribute a community pipeline](#6-contribute-a-community-pipeline) for official pipelines vs. community pipelines: It is not feasible for the core maintainers to maintain all possible training methods for diffusion models.
If the Diffusers core maintainers and the community consider a certain training paradigm to be too experimental or not popular enough, the corresponding training code should be put in the `research_projects` folder and maintained by the author.
Both official training and research examples consist of a directory that contains one or more training scripts, a requirements.txt file, and a README.md file. In order for the user to make use of the
training examples, it is required to clone the repository:
Therefore when adding an example, the `requirements.txt` file shall define all pip dependencies required for your training example so that once all those are installed, the user can run the example's training script. See, for example, the [DreamBooth `requirements.txt` file](https://github.com/huggingface/diffusers/blob/main/examples/dreambooth/requirements.txt).
Training examples of the Diffusers library should adhere to the following philosophy:
- All the code necessary to run the examples should be found in a single Python file
- One should be able to run the example from the command line with `python <your-example>.py --args`
- All the code necessary to run the examples should be found in a single Python file.
- One should be able to run the example from the command line with `python <your-example>.py --args`.
- Examples should be kept simple and serve as **an example** on how to use Diffusers for training. The purpose of example scripts is **not** to create state-of-the-art diffusion models, but rather to reproduce known training schemes without adding too much custom logic. As a byproduct of this point, our examples also strive to serve as good educational materials.
To contribute an example, it is highly recommended to look at already existing examples such as [dreambooth](https://github.com/huggingface/diffusers/blob/main/examples/dreambooth/train_dreambooth.py) to get an idea of how they should look like.
@@ -281,7 +281,7 @@ If you are contributing to the official training examples, please also make sure
usually more complicated to solve than [Good first issues](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22good+first+issue%22).
The issue description usually gives less guidance on how to fix the issue and requires
a decent understanding of the library by the interested contributor.
If you are interested in tackling a second good issue, feel free to open a PR to fix it and link the PR to the issue. If you see that a PR has already been opened for this issue but did not get merged, have a look to understand why it wasn't merged and try to open an improved PR.
If you are interested in tackling a good second issue, feel free to open a PR to fix it and link the PR to the issue. If you see that a PR has already been opened for this issue but did not get merged, have a look to understand why it wasn't merged and try to open an improved PR.
Good second issues are usually more difficult to get merged compared to good first issues, so don't hesitate to ask for help from the core maintainers. If your PR is almost finished the core maintainers can also jump into your PR and commit to it in order to get it merged.
### 9. Adding pipelines, models, schedulers
@@ -337,8 +337,8 @@ to be merged;
9. Add high-coverage tests. No quality testing = no merge.
- If you are adding new `@slow` tests, make sure they pass using
CircleCI does not run the slow tests, but GitHub actions does every night!
10. All public methods must have informative docstrings that work nicely with markdown. See `[pipeline_latent_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/latent_diffusion/pipeline_latent_diffusion.py)` for an example.
CircleCI does not run the slow tests, but GitHub Actions does every night!
10. All public methods must have informative docstrings that work nicely with markdown. See [`pipeline_latent_diffusion.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/latent_diffusion/pipeline_latent_diffusion.py) for an example.
11. Due to the rapidly growing repository, it is important to make sure that no files that would significantly weigh down the repository are added. This includes images, videos, and other non-text files. We prefer to leverage a hf.co hosted `dataset` like
[`hf-internal-testing`](https://huggingface.co/hf-internal-testing) or [huggingface/documentation-images](https://huggingface.co/datasets/huggingface/documentation-images) to place these files.
If an external contribution, feel free to add the images to your PR and ask a Hugging Face member to migrate your images
@@ -355,7 +355,7 @@ You will need basic `git` proficiency to be able to contribute to
manual. Type `git --help` in a shell and enjoy. If you prefer books, [Pro
Git](https://git-scm.com/book/en/v2) is a very good reference.
Follow these steps to start contributing ([supported Python versions](https://github.com/huggingface/diffusers/blob/main/setup.py#L244)):
Follow these steps to start contributing ([supported Python versions](https://github.com/huggingface/diffusers/blob/main/setup.py#L265)):
1. Fork the [repository](https://github.com/huggingface/diffusers) by
clicking on the 'Fork' button on the repository's page. This creates a copy of the code
@@ -364,7 +364,7 @@ under your GitHub user account.
2. Clone your fork to your local disk, and add the base repository as a remote:
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -22,7 +22,7 @@ In a nutshell, Diffusers is built to be a natural extension of PyTorch. Therefor
## Usability over Performance
- While Diffusers has many built-in performance-enhancing features (see [Memory and Speed](https://huggingface.co/docs/diffusers/optimization/fp16)), models are always loaded with the highest precision and lowest optimization. Therefore, by default diffusion pipelines are always instantiated on CPU with float32 precision if not otherwise defined by the user. This ensures usability across different platforms and accelerators and means that no complex installations are required to run the library.
- Diffusers aim at being a **light-weight** package and therefore has very few required dependencies, but many soft dependencies that can improve performance (such as `accelerate`, `safetensors`, `onnx`, etc...). We strive to keep the library as lightweight as possible so that it can be added without much concern as a dependency on other packages.
- Diffusers aims to be a **light-weight** package and therefore has very few required dependencies, but many soft dependencies that can improve performance (such as `accelerate`, `safetensors`, `onnx`, etc...). We strive to keep the library as lightweight as possible so that it can be added without much concern as a dependency on other packages.
- Diffusers prefers simple, self-explainable code over condensed, magic code. This means that short-hand code syntaxes such as lambda functions, and advanced PyTorch operators are often not desired.
## Simple over easy
@@ -31,13 +31,13 @@ As PyTorch states, **explicit is better than implicit** and **simple is better t
- We follow PyTorch's API with methods like [`DiffusionPipeline.to`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.to) to let the user handle device management.
- Raising concise error messages is preferred to silently correct erroneous input. Diffusers aims at teaching the user, rather than making the library as easy to use as possible.
- Complex model vs. scheduler logic is exposed instead of magically handled inside. Schedulers/Samplers are separated from diffusion models with minimal dependencies on each other. This forces the user to write the unrolled denoising loop. However, the separation allows for easier debugging and gives the user more control over adapting the denoising process or switching out diffusion models or schedulers.
- Separately trained components of the diffusion pipeline, *e.g.* the text encoder, the unet, and the variational autoencoder, each have their own model class. This forces the user to handle the interaction between the different model components, and the serialization format separates the model components into different files. However, this allows for easier debugging and customization. Dreambooth or textual inversion training
is very simple thanks to diffusers' ability to separate single components of the diffusion pipeline.
- Separately trained components of the diffusion pipeline, *e.g.* the text encoder, the UNet, and the variational autoencoder, each has their own model class. This forces the user to handle the interaction between the different model components, and the serialization format separates the model components into different files. However, this allows for easier debugging and customization. DreamBooth or Textual Inversion training
is very simple thanks to Diffusers' ability to separate single components of the diffusion pipeline.
## Tweakable, contributor-friendly over abstraction
For large parts of the library, Diffusers adopts an important design principle of the [Transformers library](https://github.com/huggingface/transformers), which is to prefer copy-pasted code over hasty abstractions. This design principle is very opinionated and stands in stark contrast to popular design principles such as [Don't repeat yourself (DRY)](https://en.wikipedia.org/wiki/Don%27t_repeat_yourself).
In short, just like Transformers does for modeling files, diffusers prefers to keep an extremely low level of abstraction and very self-contained code for pipelines and schedulers.
In short, just like Transformers does for modeling files, Diffusers prefers to keep an extremely low level of abstraction and very self-contained code for pipelines and schedulers.
Functions, long code blocks, and even classes can be copied across multiple files which at first can look like a bad, sloppy design choice that makes the library unmaintainable.
**However**, this design has proven to be extremely successful for Transformers and makes a lot of sense for community-driven, open-source machine learning libraries because:
- Machine Learning is an extremely fast-moving field in which paradigms, model architectures, and algorithms are changing rapidly, which therefore makes it very difficult to define long-lasting code abstractions.
@@ -47,30 +47,30 @@ Functions, long code blocks, and even classes can be copied across multiple file
At Hugging Face, we call this design the **single-file policy** which means that almost all of the code of a certain class should be written in a single, self-contained file. To read more about the philosophy, you can have a look
at [this blog post](https://huggingface.co/blog/transformers-design-philosophy).
In diffusers, we follow this philosophy for both pipelines and schedulers, but only partly for diffusion models. The reason we don't follow this design fully for diffusion models is because almost all diffusion pipelines, such
as [DDPM](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/ddpm), [Stable Diffusion](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/stable_diffusion/overview#stable-diffusion-pipelines), [UnCLIP (Dalle-2)](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/unclip#overview) and [Imagen](https://imagen.research.google/) all rely on the same diffusion model, the [UNet](https://huggingface.co/docs/diffusers/api/models#diffusers.UNet2DConditionModel).
In Diffusers, we follow this philosophy for both pipelines and schedulers, but only partly for diffusion models. The reason we don't follow this design fully for diffusion models is because almost all diffusion pipelines, such
as [DDPM](https://huggingface.co/docs/diffusers/api/pipelines/ddpm), [Stable Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/overview#stable-diffusion-pipelines), [unCLIP (DALL·E 2)](https://huggingface.co/docs/diffusers/api/pipelines/unclip) and [Imagen](https://imagen.research.google/) all rely on the same diffusion model, the [UNet](https://huggingface.co/docs/diffusers/api/models/unet2d-cond).
Great, now you should have generally understood why 🧨 Diffusers is designed the way it is 🤗.
We try to apply these design principles consistently across the library. Nevertheless, there are some minor exceptions to the philosophy or some unlucky design choices. If you have feedback regarding the design, we would ❤️ to hear it [directly on GitHub](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feedback.md&title=).
## Design Philosophy in Details
Now, let's look a bit into the nitty-gritty details of the design philosophy. Diffusers essentially consist of three major classes, [pipelines](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines), [models](https://github.com/huggingface/diffusers/tree/main/src/diffusers/models), and [schedulers](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers).
Let's walk through more in-detail design decisions for each class.
Now, let's look a bit into the nitty-gritty details of the design philosophy. Diffusers essentially consists of three major classes: [pipelines](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines), [models](https://github.com/huggingface/diffusers/tree/main/src/diffusers/models), and [schedulers](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers).
Let's walk through more detailed design decisions for each class.
### Pipelines
Pipelines are designed to be easy to use (therefore do not follow [*Simple over easy*](#simple-over-easy) 100%)), are not feature complete, and should loosely be seen as examples of how to use [models](#models) and [schedulers](#schedulers) for inference.
Pipelines are designed to be easy to use (therefore do not follow [*Simple over easy*](#simple-over-easy) 100%), are not feature complete, and should loosely be seen as examples of how to use [models](#models) and [schedulers](#schedulers) for inference.
The following design principles are followed:
- Pipelines follow the single-file policy. All pipelines can be found in individual directories under src/diffusers/pipelines. One pipeline folder corresponds to one diffusion paper/project/release. Multiple pipeline files can be gathered in one pipeline folder, as it’s done for [`src/diffusers/pipelines/stable-diffusion`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines/stable_diffusion). If pipelines share similar functionality, one can make use of the [#Copied from mechanism](https://github.com/huggingface/diffusers/blob/125d783076e5bd9785beb05367a2d2566843a271/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_img2img.py#L251).
- Pipelines all inherit from [`DiffusionPipeline`]
- Pipelines all inherit from [`DiffusionPipeline`].
- Every pipeline consists of different model and scheduler components, that are documented in the [`model_index.json` file](https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/model_index.json), are accessible under the same name as attributes of the pipeline and can be shared between pipelines with [`DiffusionPipeline.components`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.components) function.
- Every pipeline should be loadable via the [`DiffusionPipeline.from_pretrained`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.from_pretrained) function.
- Pipelines should be used **only** for inference.
- Pipelines should be very readable, self-explanatory, and easy to tweak.
- Pipelines should be designed to build on top of each other and be easy to integrate into higher-level APIs.
- Pipelines are **not** intended to be feature-complete user interfaces. For future complete user interfaces one should rather have a look at [InvokeAI](https://github.com/invoke-ai/InvokeAI), [Diffuzers](https://github.com/abhishekkrthakur/diffuzers), and [lama-cleaner](https://github.com/Sanster/lama-cleaner)
- Pipelines are **not** intended to be feature-complete user interfaces. For future complete user interfaces one should rather have a look at [InvokeAI](https://github.com/invoke-ai/InvokeAI), [Diffuzers](https://github.com/abhishekkrthakur/diffuzers), and [lama-cleaner](https://github.com/Sanster/lama-cleaner).
- Every pipeline should have one and only one way to run it via a `__call__` method. The naming of the `__call__` arguments should be shared across all pipelines.
- Pipelines should be named after the task they are intended to solve.
- In almost all cases, novel diffusion pipelines shall be implemented in a new pipeline folder/file.
@@ -83,14 +83,14 @@ The following design principles are followed:
- Models correspond to **a type of model architecture**. *E.g.* the [`UNet2DConditionModel`] class is used for all UNet variations that expect 2D image inputs and are conditioned on some context.
- All models can be found in [`src/diffusers/models`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/models) and every model architecture shall be defined in its file, e.g. [`unet_2d_condition.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_condition.py), [`transformer_2d.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/transformer_2d.py), etc...
- Models **do not** follow the single-file policy and should make use of smaller model building blocks, such as [`attention.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention.py), [`resnet.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/resnet.py), [`embeddings.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/embeddings.py), etc... **Note**: This is in stark contrast to Transformers' modeling files and shows that models do not really follow the single-file policy.
- Models intend to expose complexity, just like PyTorch's module does, and give clear error messages.
- Models intend to expose complexity, just like PyTorch's `Module` class, and give clear error messages.
- Models all inherit from `ModelMixin` and `ConfigMixin`.
- Models can be optimized for performance when it doesn’t demand major code changes, keeps backward compatibility, and gives significant memory or compute gain.
- Models can be optimized for performance when it doesn’t demand major code changes, keep backward compatibility, and give significant memory or compute gain.
- Models should by default have the highest precision and lowest performance setting.
- To integrate new model checkpoints whose general architecture can be classified as an architecture that already exists in Diffusers, the existing model architecture shall be adapted to make it work with the new checkpoint. One should only create a new file if the model architecture is fundamentally different.
- Models should be designed to be easily extendable to future changes. This can be achieved by limiting public function arguments, configuration arguments, and "foreseeing" future changes, *e.g.* it is usually better to add `string` "...type" arguments that can easily be extended to new future types instead of boolean `is_..._type` arguments. Only the minimum amount of changes shall be made to existing architectures to make a new model checkpoint work.
- The model design is a difficult trade-off between keeping code readable and concise and supporting many model checkpoints. For most parts of the modeling code, classes shall be adapted for new model checkpoints, while there are some exceptions where it is preferred to add new classes to make sure the code is kept concise and
readable longterm, such as [UNet blocks](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_blocks.py) and [Attention processors](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py).
readable long-term, such as [UNet blocks](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_blocks.py) and [Attention processors](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py).
### Schedulers
@@ -99,12 +99,12 @@ Schedulers are responsible to guide the denoising process for inference as well
The following design principles are followed:
- All schedulers are found in [`src/diffusers/schedulers`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers).
- Schedulers are **not** allowed to import from large utils files and shall be kept very self-contained.
- One scheduler python file corresponds to one scheduler algorithm (as might be defined in a paper).
- One scheduler Python file corresponds to one scheduler algorithm (as might be defined in a paper).
- If schedulers share similar functionalities, we can make use of the `#Copied from` mechanism.
- Schedulers all inherit from `SchedulerMixin` and `ConfigMixin`.
- Schedulers can be easily swapped out with the [`ConfigMixin.from_config`](https://huggingface.co/docs/diffusers/main/en/api/configuration#diffusers.ConfigMixin.from_config) method as explained in detail [here](./using-diffusers/schedulers.md).
- Schedulers can be easily swapped out with the [`ConfigMixin.from_config`](https://huggingface.co/docs/diffusers/main/en/api/configuration#diffusers.ConfigMixin.from_config) method as explained in detail [here](./docs/source/en/using-diffusers/schedulers.md).
- Every scheduler has to have a `set_num_inference_steps`, and a `step` function. `set_num_inference_steps(...)` has to be called before every denoising process, *i.e.* before `step(...)` is called.
- Every scheduler exposes the timesteps to be "looped over" via a `timesteps` attribute, which is an array of timesteps the model will be called upon
- Every scheduler exposes the timesteps to be "looped over" via a `timesteps` attribute, which is an array of timesteps the model will be called upon.
- The `step(...)` function takes a predicted model output and the "current" sample (x_t) and returns the "previous", slightly more denoised sample (x_t-1).
- Given the complexity of diffusion schedulers, the `step` function does not expose all the complexity and can be a bit of a "black box".
- In almost all cases, novel schedulers shall be implemented in a new scheduling file.
- State-of-the-art [diffusion pipelines](https://huggingface.co/docs/diffusers/api/pipelines/overview) that can be run in inference with just a few lines of code.
- Interchangeable noise [schedulers](https://huggingface.co/docs/diffusers/api/schedulers/overview) for different diffusion speeds and output quality.
- Pretrained [models](https://huggingface.co/docs/diffusers/api/models) that can be used as building blocks, and combined with schedulers, for creating your own end-to-end diffusion systems.
- Pretrained [models](https://huggingface.co/docs/diffusers/api/models/overview) that can be used as building blocks, and combined with schedulers, for creating your own end-to-end diffusion systems.
## Installation
We recommend installing 🤗 Diffusers in a virtual environment from PyPi or Conda. For more details about installing [PyTorch](https://pytorch.org/get-started/locally/) and [Flax](https://flax.readthedocs.io/en/latest/#installation), please refer to their official documentation.
We recommend installing 🤗 Diffusers in a virtual environment from PyPI or Conda. For more details about installing [PyTorch](https://pytorch.org/get-started/locally/) and [Flax](https://flax.readthedocs.io/en/latest/#installation), please refer to their official documentation.
### PyTorch
@@ -58,7 +77,7 @@ Please refer to the [How to use Stable Diffusion in Apple Silicon](https://huggi
## Quickstart
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 4000+ checkpoints):
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 19000+ checkpoints):
```python
fromdiffusersimportDiffusionPipeline
@@ -75,14 +94,13 @@ You can also dig into the models and schedulers toolbox to build your own diffus
@@ -117,8 +135,7 @@ You can look out for [issues](https://github.com/huggingface/diffusers/issues) y
- See [New model/pipeline](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22New+pipeline%2Fmodel%22) to contribute exciting new diffusion models / diffusion pipelines
- See [New scheduler](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22New+scheduler%22)
Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz98XR"><img alt="Join us on Discord" src="https://img.shields.io/discord/823813159592001537?color=5865F2&logo=discord&logoColor=white"></a>. We discuss the hottest trends about diffusion models, help each other with contributions, personal projects or
just hang out ☕.
Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz98XR"><img alt="Join us on Discord" src="https://img.shields.io/discord/823813159592001537?color=5865F2&logo=discord&logoColor=white"></a>. We discuss the hottest trends about diffusion models, help each other with contributions, personal projects or just hang out ☕.
@@ -109,8 +109,8 @@ although we can write them directly in Markdown.
Adding a new tutorial or section is done in two steps:
- Add a new file under `docs/source`. This file can either be ReStructuredText (.rst) or Markdown (.md).
- Link that file in `docs/source/_toctree.yml` on the correct toc-tree.
- Add a new Markdown (.md) file under `docs/source/<languageCode>`.
- Link that file in `docs/source/<languageCode>/_toctree.yml` on the correct toc-tree.
Make sure to put your new file under the proper section. It's unlikely to go in the first section (*Get Started*), so
depending on the intended targets (beginners, more advanced users, or researchers) it should go in sections two, three, or four.
@@ -119,7 +119,7 @@ depending on the intended targets (beginners, more advanced users, or researcher
When adding a new pipeline:
-create a file `xxx.md` under `docs/source/api/pipelines` (don't hesitate to copy an existing file as template).
-Create a file `xxx.md` under `docs/source/<languageCode>/api/pipelines` (don't hesitate to copy an existing file as template).
- Link that file in (*Diffusers Summary*) section in `docs/source/api/pipelines/overview.md`, along with the link to the paper, and a colab notebook (if available).
- Write a short overview of the diffusion model:
- Overview with paper & authors
@@ -129,8 +129,6 @@ When adding a new pipeline:
- Add all the pipeline classes that should be linked in the diffusion model. These classes should be added using our Markdown syntax. By default as follows:
```
## XXXPipeline
[[autodoc]] XXXPipeline
- all
- __call__
@@ -144,11 +142,11 @@ This will include every public method of the pipeline that is documented, as wel
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_xformers_memory_efficient_attention
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
```
You can follow the same process to create a new scheduler under the `docs/source/api/schedulers` folder
You can follow the same process to create a new scheduler under the `docs/source/<languageCode>/api/schedulers` folder.
### Writing source documentation
@@ -156,7 +154,7 @@ Values that should be put in `code` should either be surrounded by backticks: \`
and objects like True, None, or any strings should usually be put in `code`.
When mentioning a class, function, or method, it is recommended to use our syntax for internal links so that our tool
adds a link to its documentation with this syntax: \[\`XXXClass\`\] or \[\`function\`\]. This requires the class or
adds a link to its documentation with this syntax: \[\`XXXClass\`\] or \[\`function\`\]. This requires the class or
function to be in the main package.
If you want to create a link to some internal class or function, you need to
@@ -164,7 +162,7 @@ provide its path. For instance: \[\`pipelines.ImagePipelineOutput\`\]. This will
`pipelines.ImagePipelineOutput` in the description. To get rid of the path and only keep the name of the object you are
linking to in the description, add a ~: \[\`~pipelines.ImagePipelineOutput\`\] will generate a link with `ImagePipelineOutput` in the description.
The same works for methods so you can either use \[\`XXXClass.method\`\] or \[~\`XXXClass.method\`\].
The same works for methods so you can either use \[\`XXXClass.method\`\] or \[\`~XXXClass.method\`\].
#### Defining arguments in a method
@@ -196,8 +194,8 @@ Here's an example showcasing everything so far:
For optional arguments or arguments with defaults we follow the following syntax: imagine we have a function with the
following signature:
```
defmy_function(x: str = None, a: float = 1):
```py
def my_function(x:str=None,a:float=3.14):
```
then its documentation should look like this:
@@ -206,7 +204,7 @@ then its documentation should look like this:
Args:
x (`str`, *optional*):
This argument controls ...
a (`float`, *optional*, defaults to 1):
a (`float`, *optional*, defaults to `3.14`):
This argument is used to ...
```
@@ -268,4 +266,3 @@ We have an automatic script running with the `make style` command that will make
This script may have some weird failures if you made a syntax mistake or if you uncover a bug. Therefore, it's
recommended to commit your changes before running `make style`, so you can revert the changes done by that script
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the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
### Translating the Diffusers documentation into your language
As part of our mission to democratize machine learning, we'd love to make the Diffusers library available in many more languages! Follow the steps below if you want to help translate the documentation into your language 🙏.
**🗞️ Open an issue**
To get started, navigate to the [Issues](https://github.com/huggingface/diffusers/issues) page of this repo and check if anyone else has opened an issue for your language. If not, open a new issue by selecting the "Translation template" from the "New issue" button.
To get started, navigate to the [Issues](https://github.com/huggingface/diffusers/issues) page of this repo and check if anyone else has opened an issue for your language. If not, open a new issue by selecting the "🌐 Translating a New Language?" from the "New issue" button.
Once an issue exists, post a comment to indicate which chapters you'd like to work on, and we'll add your name to the list.
@@ -16,7 +28,7 @@ First, you'll need to [fork the Diffusers repo](https://docs.github.com/en/get-s
Once you've forked the repo, you'll want to get the files on your local machine for editing. You can do that by cloning the fork with Git as follows:
**📋 Copy-paste the English version with a new language code**
@@ -29,18 +41,18 @@ You'll only need to copy the files in the [`docs/source/en`](https://github.com/
```bash
cd ~/path/to/diffusers/docs
cp -r source/en source/LANG-ID
cp -r source/en source/<LANG-ID>
```
Here, `LANG-ID` should be one of the ISO 639-1 or ISO 639-2 language codes -- see [here](https://www.loc.gov/standards/iso639-2/php/code_list.php) for a handy table.
Here, `<LANG-ID>` should be one of the ISO 639-1 or ISO 639-2 language codes -- see [here](https://www.loc.gov/standards/iso639-2/php/code_list.php) for a handy table.
**✍️ Start translating**
The fun part comes - translating the text!
The first thing we recommend is translating the part of the `_toctree.yml` file that corresponds to your doc chapter. This file is used to render the table of contents on the website.
The first thing we recommend is translating the part of the `_toctree.yml` file that corresponds to your doc chapter. This file is used to render the table of contents on the website.
> 🙋 If the `_toctree.yml` file doesn't yet exist for your language, you can create one by copy-pasting from the English version and deleting the sections unrelated to your chapter. Just make sure it exists in the `docs/source/LANG-ID/` directory!
> 🙋 If the `_toctree.yml` file doesn't yet exist for your language, you can create one by copy-pasting from the English version and deleting the sections unrelated to your chapter. Just make sure it exists in the `docs/source/<LANG-ID>/` directory!
The fields you should add are `local` (with the name of the file containing the translation; e.g. `autoclass_tutorial`), and `title` (with the title of the doc in your language; e.g. `Load pretrained instances with an AutoClass`) -- as a reference, here is the `_toctree.yml` for [English](https://github.com/huggingface/diffusers/blob/main/docs/source/en/_toctree.yml):
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Pipelines
The [`DiffusionPipeline`] is the quickest way to load any pretrained diffusion pipeline from the [Hub](https://huggingface.co/models?library=diffusers) for inference.
<Tip>
You shouldn't use the [`DiffusionPipeline`] class for training or finetuning a diffusion model. Individual
components (for example, [`UNet2DModel`] and [`UNet2DConditionModel`]) of diffusion pipelines are usually trained individually, so we suggest directly working with them instead.
</Tip>
The pipeline type (for example [`StableDiffusionPipeline`]) of any diffusion pipeline loaded with [`~DiffusionPipeline.from_pretrained`] is automatically
detected and pipeline components are loaded and passed to the `__init__` function of the pipeline.
Any pipeline object can be saved locally with [`~DiffusionPipeline.save_pretrained`].
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,9 +12,9 @@ specific language governing permissions and limitations under the License.
# VAE Image Processor
The [`VaeImageProcessor`] provides a unified API for [`StableDiffusionPipeline`]'s to prepare image inputs for VAE encoding and post-processing outputs once they're decoded. This includes transformations such as resizing, normalization, and conversion between PIL Image, PyTorch, and NumPy arrays.
The [`VaeImageProcessor`] provides a unified API for [`StableDiffusionPipeline`]s to prepare image inputs for VAE encoding and post-processing outputs once they're decoded. This includes transformations such as resizing, normalization, and conversion between PIL Image, PyTorch, and NumPy arrays.
All pipelines with [`VaeImageProcessor`] accepts PIL Image, PyTorch tensor, or NumPy arrays as image inputs and returns outputs based on the `output_type` argument by the user. You can pass encoded image latents directly to the pipeline and return latents from the pipeline as a specific output with the `output_type` argument (for example `output_type="pt"`). This allows you to take the generated latents from one pipeline and pass it to another pipeline as input without leaving the latent space. It also makes it much easier to use multiple pipelines together by passing PyTorch tensors directly between different pipelines.
All pipelines with [`VaeImageProcessor`] accept PIL Image, PyTorch tensor, or NumPy arrays as image inputs and return outputs based on the `output_type` argument by the user. You can pass encoded image latents directly to the pipeline and return latents from the pipeline as a specific output with the `output_type` argument (for example `output_type="latent"`). This allows you to take the generated latents from one pipeline and pass it to another pipeline as input without leaving the latent space. It also makes it much easier to use multiple pipelines together by passing PyTorch tensors directly between different pipelines.
## VaeImageProcessor
@@ -24,4 +24,4 @@ All pipelines with [`VaeImageProcessor`] accepts PIL Image, PyTorch tensor, or N
The [`VaeImageProcessorLDM3D`] accepts RGB and depth inputs and returns RGB and depth outputs.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Overview
The APIs in this section are more experimental and prone to breaking changes. Most of them are used internally for development, but they may also be useful to you if you're interested in building a diffusion model with some custom parts or if you're interested in some of our helper utilities for working with 🤗 Diffusers.
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the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Loaders
Adapters (textual inversion, LoRA, hypernetworks) allow you to modify a diffusion model to generate images in a specific style without training or finetuning the entire model. The adapter weights are typically only a tiny fraction of the pretrained model's which making them very portable. 🤗 Diffusers provides an easy-to-use `LoaderMixin` API to load adapter weights.
<Tip warning={true}>
🧪 The `LoaderMixins` are highly experimental and prone to future changes. To use private or [gated](https://huggingface.co/docs/hub/models-gated#gated-models) models, log-in with `huggingface-cli login`.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# IP-Adapter
[IP-Adapter](https://hf.co/papers/2308.06721) is a lightweight adapter that enables prompting a diffusion model with an image. This method decouples the cross-attention layers of the image and text features. The image features are generated from an image encoder.
<Tip>
Learn how to load an IP-Adapter checkpoint and image in the IP-Adapter [loading](../../using-diffusers/loading_adapters#ip-adapter) guide, and you can see how to use it in the [usage](../../using-diffusers/ip_adapter) guide.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# LoRA
LoRA is a fast and lightweight training method that inserts and trains a significantly smaller number of parameters instead of all the model parameters. This produces a smaller file (~100 MBs) and makes it easier to quickly train a model to learn a new concept. LoRA weights are typically loaded into the UNet, text encoder or both. There are two classes for loading LoRA weights:
- [`LoraLoaderMixin`] provides functions for loading and unloading, fusing and unfusing, enabling and disabling, and more functions for managing LoRA weights. This class can be used with any model.
- [`StableDiffusionXLLoraLoaderMixin`] is a [Stable Diffusion (SDXL)](../../api/pipelines/stable_diffusion/stable_diffusion_xl) version of the [`LoraLoaderMixin`] class for loading and saving LoRA weights. It can only be used with the SDXL model.
<Tip>
To learn more about how to load LoRA weights, see the [LoRA](../../using-diffusers/loading_adapters#lora) loading guide.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# PEFT
Diffusers supports loading adapters such as [LoRA](../../using-diffusers/loading_adapters) with the [PEFT](https://huggingface.co/docs/peft/index) library with the [`~loaders.peft.PeftAdapterMixin`] class. This allows modeling classes in Diffusers like [`UNet2DConditionModel`] to load an adapter.
<Tip>
Refer to the [Inference with PEFT](../../tutorials/using_peft_for_inference.md) tutorial for an overview of how to use PEFT in Diffusers for inference.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Single files
Diffusers supports loading pretrained pipeline (or model) weights stored in a single file, such as a `ckpt` or `safetensors` file. These single file types are typically produced from community trained models. There are three classes for loading single file weights:
- [`FromSingleFileMixin`] supports loading pretrained pipeline weights stored in a single file, which can either be a `ckpt` or `safetensors` file.
- [`FromOriginalVAEMixin`] supports loading a pretrained [`AutoencoderKL`] from pretrained ControlNet weights stored in a single file, which can either be a `ckpt` or `safetensors` file.
- [`FromOriginalControlnetMixin`] supports loading pretrained ControlNet weights stored in a single file, which can either be a `ckpt` or `safetensors` file.
<Tip>
To learn more about how to load single file weights, see the [Load different Stable Diffusion formats](../../using-diffusers/other-formats) loading guide.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Textual Inversion
Textual Inversion is a training method for personalizing models by learning new text embeddings from a few example images. The file produced from training is extremely small (a few KBs) and the new embeddings can be loaded into the text encoder.
[`TextualInversionLoaderMixin`] provides a function for loading Textual Inversion embeddings from Diffusers and Automatic1111 into the text encoder and loading a special token to activate the embeddings.
<Tip>
To learn more about how to load Textual Inversion embeddings, see the [Textual Inversion](../../using-diffusers/loading_adapters#textual-inversion) loading guide.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# UNet
Some training methods - like LoRA and Custom Diffusion - typically target the UNet's attention layers, but these training methods can also target other non-attention layers. Instead of training all of a model's parameters, only a subset of the parameters are trained, which is faster and more efficient. This class is useful if you're *only* loading weights into a UNet. If you need to load weights into the text encoder or a text encoder and UNet, try using the [`~loaders.LoraLoaderMixin.load_lora_weights`] function instead.
The [`UNet2DConditionLoadersMixin`] class provides functions for loading and saving weights, fusing and unfusing LoRAs, disabling and enabling LoRAs, and setting and deleting adapters.
<Tip>
To learn more about how to load LoRA weights, see the [LoRA](../../using-diffusers/loading_adapters#lora) loading guide.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# AsymmetricAutoencoderKL
Improved larger variational autoencoder (VAE) model with KL loss for inpainting task: [Designing a Better Asymmetric VQGAN for StableDiffusion](https://arxiv.org/abs/2306.04632) by Zixin Zhu, Xuelu Feng, Dongdong Chen, Jianmin Bao, Le Wang, Yinpeng Chen, Lu Yuan, Gang Hua.
@@ -6,7 +18,7 @@ The abstract from the paper is:
*StableDiffusion is a revolutionary text-to-image generator that is causing a stir in the world of image generation and editing. Unlike traditional methods that learn a diffusion model in pixel space, StableDiffusion learns a diffusion model in the latent space via a VQGAN, ensuring both efficiency and quality. It not only supports image generation tasks, but also enables image editing for real images, such as image inpainting and local editing. However, we have observed that the vanilla VQGAN used in StableDiffusion leads to significant information loss, causing distortion artifacts even in non-edited image regions. To this end, we propose a new asymmetric VQGAN with two simple designs. Firstly, in addition to the input from the encoder, the decoder contains a conditional branch that incorporates information from task-specific priors, such as the unmasked image region in inpainting. Secondly, the decoder is much heavier than the encoder, allowing for more detailed recovery while only slightly increasing the total inference cost. The training cost of our asymmetric VQGAN is cheap, and we only need to retrain a new asymmetric decoder while keeping the vanilla VQGAN encoder and StableDiffusion unchanged. Our asymmetric VQGAN can be widely used in StableDiffusion-based inpainting and local editing methods. Extensive experiments demonstrate that it can significantly improve the inpainting and editing performance, while maintaining the original text-to-image capability. The code is available at https://github.com/buxiangzhiren/Asymmetric_VQGAN*
Evaluation results can be found in section 4.1 of the original paper.
Evaluation results can be found in section 4.1 of the original paper.
## Available checkpoints
@@ -16,40 +28,33 @@ Evaluation results can be found in section 4.1 of the original paper.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Tiny AutoEncoder
Tiny AutoEncoder for Stable Diffusion (TAESD) was introduced in [madebyollin/taesd](https://github.com/madebyollin/taesd) by Ollin Boer Bohan. It is a tiny distilled version of Stable Diffusion's VAE that can quickly decode the latents in a [`StableDiffusionPipeline`] or [`StableDiffusionXLPipeline`] almost instantly.
Tiny AutoEncoder for Stable Diffusion (TAESD) was introduced in [madebyollin/taesd](https://github.com/madebyollin/taesd) by Ollin Boer Bohan. It is a tiny distilled version of Stable Diffusion's VAE that can quickly decode the latents in a [`StableDiffusionPipeline`] or [`StableDiffusionXLPipeline`] almost instantly.
To use with Stable Diffusion v-2.1:
@@ -16,7 +28,7 @@ pipe = pipe.to("cuda")
prompt="slice of delicious New York-style berry cheesecake"
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# AutoencoderKL
The variational autoencoder (VAE) model with KL loss was introduced in [Auto-Encoding Variational Bayes](https://arxiv.org/abs/1312.6114v11) by Diederik P. Kingma and Max Welling. The model is used in 🤗 Diffusers to encode images into latents and to decode latent representations into images.
@@ -14,21 +26,24 @@ from the original format using [`FromOriginalVAEMixin.from_single_file`] as foll
```py
fromdiffusersimportAutoencoderKL
url="https://huggingface.co/stabilityai/sd-vae-ft-mse-original/blob/main/vae-ft-mse-840000-ema-pruned.safetensors"# can also be local file
url="https://huggingface.co/stabilityai/sd-vae-ft-mse-original/blob/main/vae-ft-mse-840000-ema-pruned.safetensors"# can also be a local file
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Consistency Decoder
Consistency decoder can be used to decode the latents from the denoising UNet in the [`StableDiffusionPipeline`]. This decoder was introduced in the [DALL-E 3 technical report](https://openai.com/dall-e-3).
The original codebase can be found at [openai/consistencydecoder](https://github.com/openai/consistencydecoder).
<Tip warning={true}>
Inference is only supported for 2 iterations as of now.
</Tip>
The pipeline could not have been contributed without the help of [madebyollin](https://github.com/madebyollin) and [mrsteyk](https://github.com/mrsteyk) from [this issue](https://github.com/openai/consistencydecoder/issues/1).
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
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specific language governing permissions and limitations under the License.
-->
# ControlNet
The ControlNet model was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang and Maneesh Agrawala. It provides a greater degree of control over text-to-image generation by conditioning the model on additional inputs such as edge maps, depth maps, segmentation maps, and keypoints for pose detection.
The ControlNet model was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, Maneesh Agrawala. It provides a greater degree of control over text-to-image generation by conditioning the model on additional inputs such as edge maps, depth maps, segmentation maps, and keypoints for pose detection.
The abstract from the paper is:
*We present a neural network structure, ControlNet, to control pretrained large diffusion models to support additional input conditions. The ControlNet learns task-specific conditions in an end-to-end way, and the learning is robust even when the training dataset is small (< 50k). Moreover, training a ControlNet is as fast as fine-tuning a diffusion model, and the model can be trained on a personal devices. Alternatively, if powerful computation clusters are available, the model can scale to large amounts (millions to billions) of data. We report that large diffusion models like Stable Diffusion can be augmented with ControlNets to enable conditional inputs like edge maps, segmentation maps, keypoints, etc. This may enrich the methods to control large diffusion models and further facilitate related applications.*
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
## Loading from the original format
@@ -12,13 +24,13 @@ By default the [`ControlNetModel`] should be loaded with [`~ModelMixin.from_pret
from the original format using [`FromOriginalControlnetMixin.from_single_file`] as follows:
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# Models
🤗 Diffusers provides pretrained models for popular algorithms and modules to create custom diffusion systems. The primary function of models is to denoise an input sample as modeled by the distribution \\(p_{\theta}(x_{t-1}|x_{t})\\).
🤗 Diffusers provides pretrained models for popular algorithms and modules to create custom diffusion systems. The primary function of models is to denoise an input sample as modeled by the distribution \\(p_{\theta}(x_{t-1}|x_{t})\\).
All models are built from the base [`ModelMixin`] class which is a [`torch.nn.module`](https://pytorch.org/docs/stable/generated/torch.nn.Module.html) providing basic functionality for saving and loading models, locally and from the Hugging Face Hub.
All models are built from the base [`ModelMixin`] class which is a [`torch.nn.Module`](https://pytorch.org/docs/stable/generated/torch.nn.Module.html) providing basic functionality for saving and loading models, locally and from the Hugging Face Hub.
## ModelMixin
[[autodoc]] ModelMixin
@@ -13,4 +25,4 @@ All models are built from the base [`ModelMixin`] class which is a [`torch.nn.mo
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# Prior Transformer
The Prior Transformer was originally introduced in [Hierarchical Text-Conditional Image Generation with CLIP Latents
](https://huggingface.co/papers/2204.06125) by Ramesh et al. It is used to predict CLIP image embeddings from CLIP text embeddings; image embeddings are predicted through a denoising diffusion process.
The Prior Transformer was originally introduced in [Hierarchical Text-Conditional Image Generation with CLIP Latents](https://huggingface.co/papers/2204.06125) by Ramesh et al. It is used to predict CLIP image embeddings from CLIP text embeddings; image embeddings are predicted through a denoising diffusion process.
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# Transformer2D
A Transformer model for image-like data from [CompVis](https://huggingface.co/CompVis) that is based on the [Vision Transformer](https://huggingface.co/papers/2010.11929) introduced by Dosovitskiy et al. The [`Transformer2DModel`] accepts discrete (classes of vector embeddings) or continuous (actual embeddings) inputs.
@@ -26,4 +38,4 @@ It is assumed one of the input classes is the masked latent pixel. The predicted
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# UNetMotionModel
The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 2D UNet model.
The abstract from the paper is:
*There is large consent that successful training of deep networks requires many thousand annotated training samples. In this paper, we present a network and training strategy that relies on the strong use of data augmentation to use the available annotated samples more efficiently. The architecture consists of a contracting path to capture context and a symmetric expanding path that enables precise localization. We show that such a network can be trained end-to-end from very few images and outperforms the prior best method (a sliding-window convolutional network) on the ISBI challenge for segmentation of neuronal structures in electron microscopic stacks. Using the same network trained on transmitted light microscopy images (phase contrast and DIC) we won the ISBI cell tracking challenge 2015 in these categories by a large margin. Moreover, the network is fast. Segmentation of a 512x512 image takes less than a second on a recent GPU. The full implementation (based on Caffe) and the trained networks are available at http://lmb.informatik.uni-freiburg.de/people/ronneber/u-net.*
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# UNet1DModel
The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 1D UNet model.
The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al. for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 1D UNet model.
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# UNet2DConditionModel
The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 2D UNet conditional model.
The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al. for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 2D UNet conditional model.
The abstract from the paper is:
@@ -10,10 +22,10 @@ The abstract from the paper is:
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# UNet2DModel
The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 2D UNet model.
The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al. for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 2D UNet model.
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# UNet3DConditionModel
The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 3D UNet conditional model.
The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al. for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 3D UNet conditional model.
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# UVit2DModel
The [U-ViT](https://hf.co/papers/2301.11093) model is a vision transformer (ViT) based UNet. This model incorporates elements from ViT (considers all inputs such as time, conditions and noisy image patches as tokens) and a UNet (long skip connections between the shallow and deep layers). The skip connection is important for predicting pixel-level features. An additional 3x3 convolutional block is applied prior to the final output to improve image quality.
The abstract from the paper is:
*Currently, applying diffusion models in pixel space of high resolution images is difficult. Instead, existing approaches focus on diffusion in lower dimensional spaces (latent diffusion), or have multiple super-resolution levels of generation referred to as cascades. The downside is that these approaches add additional complexity to the diffusion framework. This paper aims to improve denoising diffusion for high resolution images while keeping the model as simple as possible. The paper is centered around the research question: How can one train a standard denoising diffusion models on high resolution images, and still obtain performance comparable to these alternate approaches? The four main findings are: 1) the noise schedule should be adjusted for high resolution images, 2) It is sufficient to scale only a particular part of the architecture, 3) dropout should be added at specific locations in the architecture, and 4) downsampling is an effective strategy to avoid high resolution feature maps. Combining these simple yet effective techniques, we achieve state-of-the-art on image generation among diffusion models without sampling modifiers on ImageNet.*
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# VQModel
The VQ-VAE model was introduced in [Neural Discrete Representation Learning](https://huggingface.co/papers/1711.00937) by Aaron van den Oord, Oriol Vinyals and Koray Kavukcuoglu. The model is used in 🤗 Diffusers to decode latent representations into images. Unlike [`AutoencoderKL`], the [`VQModel`] works in a quantized latent space.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,7 +12,7 @@ specific language governing permissions and limitations under the License.
# Outputs
All models outputs are subclasses of [`~utils.BaseOutput`], data structures containing all the information returned by the model. The outputs can also be used as tuples or dictionaries.
All model outputs are subclasses of [`~utils.BaseOutput`], data structures containing all the information returned by the model. The outputs can also be used as tuples or dictionaries.
For example:
@@ -64,4 +64,4 @@ To check a specific pipeline or model output, refer to its corresponding API doc
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# AltDiffusion
AltDiffusion was proposed in [AltCLIP: Altering the Language Encoder in CLIP for Extended Language Capabilities](https://huggingface.co/papers/2211.06679) by Zhongzhi Chen, Guang Liu, Bo-Wen Zhang, Fulong Ye, Qinghong Yang, Ledell Wu.
The abstract from the paper is:
*In this work, we present a conceptually simple and effective method to train a strong bilingual multimodal representation model. Starting from the pretrained multimodal representation model CLIP released by OpenAI, we switched its text encoder with a pretrained multilingual text encoder XLM-R, and aligned both languages and image representations by a two-stage training schema consisting of teacher learning and contrastive learning. We validate our method through evaluations of a wide range of tasks. We set new state-of-the-art performances on a bunch of tasks including ImageNet-CN, Flicker30k- CN, and COCO-CN. Further, we obtain very close performances with CLIP on almost all tasks, suggesting that one can simply alter the text encoder in CLIP for extended capabilities such as multilingual understanding.*
## Tips
`AltDiffusion` is conceptually the same as [Stable Diffusion](./stable_diffusion/overview).
<Tip>
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
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# aMUSEd
aMUSEd was introduced in [aMUSEd: An Open MUSE Reproduction](https://huggingface.co/papers/2401.01808) by Suraj Patil, William Berman, Robin Rombach, and Patrick von Platen.
Amused is a lightweight text to image model based off of the [MUSE](https://arxiv.org/abs/2301.00704) architecture. Amused is particularly useful in applications that require a lightweight and fast model such as generating many images quickly at once.
Amused is a vqvae token based transformer that can generate an image in fewer forward passes than many diffusion models. In contrast with muse, it uses the smaller text encoder CLIP-L/14 instead of t5-xxl. Due to its small parameter count and few forward pass generation process, amused can generate many images quickly. This benefit is seen particularly at larger batch sizes.
The abstract from the paper is:
*We present aMUSEd, an open-source, lightweight masked image model (MIM) for text-to-image generation based on MUSE. With 10 percent of MUSE's parameters, aMUSEd is focused on fast image generation. We believe MIM is under-explored compared to latent diffusion, the prevailing approach for text-to-image generation. Compared to latent diffusion, MIM requires fewer inference steps and is more interpretable. Additionally, MIM can be fine-tuned to learn additional styles with only a single image. We hope to encourage further exploration of MIM by demonstrating its effectiveness on large-scale text-to-image generation and releasing reproducible training code. We also release checkpoints for two models which directly produce images at 256x256 and 512x512 resolutions.*
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# Text-to-Video Generation with AnimateDiff
## Overview
[AnimateDiff: Animate Your Personalized Text-to-Image Diffusion Models without Specific Tuning](https://arxiv.org/abs/2307.04725) by Yuwei Guo, Ceyuan Yang, Anyi Rao, Yaohui Wang, Yu Qiao, Dahua Lin, Bo Dai.
The abstract of the paper is the following:
*With the advance of text-to-image models (e.g., Stable Diffusion) and corresponding personalization techniques such as DreamBooth and LoRA, everyone can manifest their imagination into high-quality images at an affordable cost. Subsequently, there is a great demand for image animation techniques to further combine generated static images with motion dynamics. In this report, we propose a practical framework to animate most of the existing personalized text-to-image models once and for all, saving efforts in model-specific tuning. At the core of the proposed framework is to insert a newly initialized motion modeling module into the frozen text-to-image model and train it on video clips to distill reasonable motion priors. Once trained, by simply injecting this motion modeling module, all personalized versions derived from the same base T2I readily become text-driven models that produce diverse and personalized animated images. We conduct our evaluation on several public representative personalized text-to-image models across anime pictures and realistic photographs, and demonstrate that our proposed framework helps these models generate temporally smooth animation clips while preserving the domain and diversity of their outputs. Code and pre-trained weights will be publicly available at [this https URL](https://animatediff.github.io/).*
## Available Pipelines
| Pipeline | Tasks | Demo
|---|---|:---:|
| [AnimateDiffPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff.py) | *Text-to-Video Generation with AnimateDiff* |
| [AnimateDiffVideoToVideoPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py) | *Video-to-Video Generation with AnimateDiff* |
## Available checkpoints
Motion Adapter checkpoints can be found under [guoyww](https://huggingface.co/guoyww/). These checkpoints are meant to work with any model based on Stable Diffusion 1.4/1.5.
## Usage example
### AnimateDiffPipeline
AnimateDiff works with a MotionAdapter checkpoint and a Stable Diffusion model checkpoint. The MotionAdapter is a collection of Motion Modules that are responsible for adding coherent motion across image frames. These modules are applied after the Resnet and Attention blocks in Stable Diffusion UNet.
The following example demonstrates how to use a *MotionAdapter* checkpoint with Diffusers for inference based on StableDiffusion-1.4/1.5.
AnimateDiff tends to work better with finetuned Stable Diffusion models. If you plan on using a scheduler that can clip samples, make sure to disable it by setting `clip_sample=False` in the scheduler as this can also have an adverse effect on generated samples. Additionally, the AnimateDiff checkpoints can be sensitive to the beta schedule of the scheduler. We recommend setting this to `linear`.
</Tip>
### AnimateDiffVideoToVideoPipeline
AnimateDiff can also be used to generate visually similar videos or enable style/character/background or other edits starting from an initial video, allowing you to seamlessly explore creative possibilities.
alt="closeup of tony stark, robert downey jr, fireworks"
style="width: 300px;" />
</td>
</tr>
</table>
## Using Motion LoRAs
Motion LoRAs are a collection of LoRAs that work with the `guoyww/animatediff-motion-adapter-v1-5-2` checkpoint. These LoRAs are responsible for adding specific types of motion to the animations.
[FreeInit: Bridging Initialization Gap in Video Diffusion Models](https://arxiv.org/abs/2312.07537) by Tianxing Wu, Chenyang Si, Yuming Jiang, Ziqi Huang, Ziwei Liu.
FreeInit is an effective method that improves temporal consistency and overall quality of videos generated using video-diffusion-models without any addition training. It can be applied to AnimateDiff, ModelScope, VideoCrafter and various other video generation models seamlessly at inference time, and works by iteratively refining the latent-initialization noise. More details can be found it the paper.
The following example demonstrates the usage of FreeInit.
FreeInit is not really free - the improved quality comes at the cost of extra computation. It requires sampling a few extra times depending on the `num_iters` parameter that is set when enabling it. Setting the `use_fast_sampling` parameter to `True` can improve the overall performance (at the cost of lower quality compared to when `use_fast_sampling=False` but still better results than vanilla video generation models).
</Tip>
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Using AnimateLCM
[AnimateLCM](https://animatelcm.github.io/) is a motion module checkpoint and an [LCM LoRA](https://huggingface.co/docs/diffusers/using-diffusers/inference_with_lcm_lora) that have been created using a consistency learning strategy that decouples the distillation of the image generation priors and the motion generation priors.
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@@ -16,13 +16,13 @@ Attend-and-Excite for Stable Diffusion was proposed in [Attend-and-Excite: Atten
The abstract from the paper is:
*Text-to-image diffusion models have recently received a lot of interest for their astonishing ability to produce high-fidelity images from text only. However, achieving one-shot generation that aligns with the user's intent is nearly impossible, yet small changes to the input prompt often result in very different images. This leaves the user with little semantic control. To put the user in control, we show how to interact with the diffusion process to flexibly steer it along semantic directions. This semantic guidance (SEGA) allows for subtle and extensive edits, changes in composition and style, as well as optimizing the overall artistic conception. We demonstrate SEGA's effectiveness on a variety of tasks and provide evidence for its versatility and flexibility.*
*Recent text-to-image generative models have demonstrated an unparalleled ability to generate diverse and creative imagery guided by a target text prompt. While revolutionary, current state-of-the-art diffusion models may still fail in generating images that fully convey the semantics in the given text prompt. We analyze the publicly available Stable Diffusion model and assess the existence of catastrophic neglect, where the model fails to generate one or more of the subjects from the input prompt. Moreover, we find that in some cases the model also fails to correctly bind attributes (e.g., colors) to their corresponding subjects. To help mitigate these failure cases, we introduce the concept of Generative Semantic Nursing (GSN), where we seek to intervene in the generative process on the fly during inference time to improve the faithfulness of the generated images. Using an attention-based formulation of GSN, dubbed Attend-and-Excite, we guide the model to refine the cross-attention units to attend to all subject tokens in the text prompt and strengthen - or excite - their activations, encouraging the model to generate all subjects described in the text prompt. We compare our approach to alternative approaches and demonstrate that it conveys the desired concepts more faithfully across a range of text prompts.*
You can find additional information about Attend-and-Excite on the [project page](https://attendandexcite.github.io/Attend-and-Excite/), the [original codebase](https://github.com/AttendAndExcite/Attend-and-Excite), or try it out in a [demo](https://huggingface.co/spaces/AttendAndExcite/Attend-and-Excite).
<Tip>
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
@@ -34,4 +34,4 @@ Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to le
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# Audio Diffusion
[Audio Diffusion](https://github.com/teticio/audio-diffusion) is by Robert Dargavel Smith, and it leverages the recent advances in image generation from diffusion models by converting audio samples to and from Mel spectrogram images.
The original codebase, training scripts and example notebooks can be found at [teticio/audio-diffusion](https://github.com/teticio/audio-diffusion).
<Tip>
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -19,9 +19,9 @@ sound effects, human speech and music.
The abstract from the paper is:
*Text-to-audio (TTA) system has recently gained attention for its ability to synthesize general audio based on text descriptions. However, previous studies in TTA have limited generation quality with high computational costs. In this study, we propose AudioLDM, a TTA system that is built on a latent space to learn the continuous audio representations from contrastive language-audio pretraining (CLAP) latents. The pretrained CLAP models enable us to train LDMs with audio embedding while providing text embedding as a condition during sampling. By learning the latent representations of audio signals and their compositions without modeling the cross-modal relationship, AudioLDM is advantageous in both generation quality and computational efficiency. Trained on AudioCaps with a single GPU, AudioLDM achieves state-of-the-art TTA performance measured by both objective and subjective metrics (e.g., frechet distance). Moreover, AudioLDM is the first TTA system that enables various text-guided audio manipulations (e.g., style transfer) in a zero-shot fashion. Our implementation and demos are available at https://audioldm.github.io.*
*Text-to-audio (TTA) system has recently gained attention for its ability to synthesize general audio based on text descriptions. However, previous studies in TTA have limited generation quality with high computational costs. In this study, we propose AudioLDM, a TTA system that is built on a latent space to learn the continuous audio representations from contrastive language-audio pretraining (CLAP) latents. The pretrained CLAP models enable us to train LDMs with audio embedding while providing text embedding as a condition during sampling. By learning the latent representations of audio signals and their compositions without modeling the cross-modal relationship, AudioLDM is advantageous in both generation quality and computational efficiency. Trained on AudioCaps with a single GPU, AudioLDM achieves state-of-the-art TTA performance measured by both objective and subjective metrics (e.g., frechet distance). Moreover, AudioLDM is the first TTA system that enables various text-guided audio manipulations (e.g., style transfer) in a zero-shot fashion. Our implementation and demos are available at [this https URL](https://audioldm.github.io/).*
The original codebase can be found at [haoheliu/AudioLDM](https://github.com/haoheliu/AudioLDM).
The original codebase can be found at [haoheliu/AudioLDM](https://github.com/haoheliu/AudioLDM).
## Tips
@@ -37,7 +37,7 @@ During inference:
<Tip>
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
@@ -47,4 +47,4 @@ Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to le
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,36 +12,23 @@ specific language governing permissions and limitations under the License.
# AudioLDM 2
AudioLDM 2 was proposed in [AudioLDM 2: Learning Holistic Audio Generation with Self-supervised Pretraining](https://arxiv.org/abs/2308.05734)
by Haohe Liu et al. AudioLDM 2 takes a text prompt as input and predicts the corresponding audio. It can generate
text-conditional sound effects, human speech and music.
AudioLDM 2 was proposed in [AudioLDM 2: Learning Holistic Audio Generation with Self-supervised Pretraining](https://arxiv.org/abs/2308.05734) by Haohe Liu et al. AudioLDM 2 takes a text prompt as input and predicts the corresponding audio. It can generate text-conditional sound effects, human speech and music.
Inspired by [Stable Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/overview), AudioLDM 2
is a text-to-audio _latent diffusion model (LDM)_ that learns continuous audio representations from text embeddings. Two
text encoder models are used to compute the text embeddings from a prompt input: the text-branch of [CLAP](https://huggingface.co/docs/transformers/main/en/model_doc/clap)
and the encoder of [Flan-T5](https://huggingface.co/docs/transformers/main/en/model_doc/flan-t5). These text embeddings
are then projected to a shared embedding space by an [AudioLDM2ProjectionModel](https://huggingface.co/docs/diffusers/main/api/pipelines/audioldm2#diffusers.AudioLDM2ProjectionModel).
A [GPT2](https://huggingface.co/docs/transformers/main/en/model_doc/gpt2) _language model (LM)_ is used to auto-regressively
predict eight new embedding vectors, conditional on the projected CLAP and Flan-T5 embeddings. The generated embedding
vectors and Flan-T5 text embeddings are used as cross-attention conditioning in the LDM. The [UNet](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2UNet2DConditionModel)
of AudioLDM 2 is unique in the sense that it takes **two** cross-attention embeddings, as opposed to one cross-attention
conditioning, as in most other LDMs.
Inspired by [Stable Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/overview), AudioLDM 2 is a text-to-audio _latent diffusion model (LDM)_ that learns continuous audio representations from text embeddings. Two text encoder models are used to compute the text embeddings from a prompt input: the text-branch of [CLAP](https://huggingface.co/docs/transformers/main/en/model_doc/clap) and the encoder of [Flan-T5](https://huggingface.co/docs/transformers/main/en/model_doc/flan-t5). These text embeddings are then projected to a shared embedding space by an [AudioLDM2ProjectionModel](https://huggingface.co/docs/diffusers/main/api/pipelines/audioldm2#diffusers.AudioLDM2ProjectionModel). A [GPT2](https://huggingface.co/docs/transformers/main/en/model_doc/gpt2) _language model (LM)_ is used to auto-regressively predict eight new embedding vectors, conditional on the projected CLAP and Flan-T5 embeddings. The generated embedding vectors and Flan-T5 text embeddings are used as cross-attention conditioning in the LDM. The [UNet](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2UNet2DConditionModel) of AudioLDM 2 is unique in the sense that it takes **two** cross-attention embeddings, as opposed to one cross-attention conditioning, as in most other LDMs.
The abstract of the paper is the following:
*Although audio generation shares commonalities across different types of audio, such as speech, music, and sound effects, designing models for each type requires careful consideration of specific objectives and biases that can significantly differ from those of other types. To bring us closer to a unified perspective of audio generation, this paper proposes a framework that utilizes the same learning method for speech, music, and sound effect generation. Our framework introduces a general representation of audio, called language of audio (LOA). Any audio can be translated into LOA based on AudioMAE, a self-supervised pre-trained representation learning model. In the generation process, we translate any modalities into LOA by using a GPT-2 model, and we perform self-supervised audio generation learning with a latent diffusion model conditioned on LOA. The proposed framework naturally brings advantages such as in-context learning abilities and reusable self-supervised pretrained AudioMAE and latent diffusion models. Experiments on the major benchmarks of text-to-audio, text-to-music, and text-to-speech demonstrate new state-of-the-art or competitive performance to previous approaches.*
*Although audio generation shares commonalities across different types of audio, such as speech, music, and sound effects, designing models for each type requires careful consideration of specific objectives and biases that can significantly differ from those of other types. To bring us closer to a unified perspective of audio generation, this paper proposes a framework that utilizes the same learning method for speech, music, and sound effect generation. Our framework introduces a general representation of audio, called "language of audio" (LOA). Any audio can be translated into LOA based on AudioMAE, a self-supervised pre-trained representation learning model. In the generation process, we translate any modalities into LOA by using a GPT-2 model, and we perform self-supervised audio generation learning with a latent diffusion model conditioned on LOA. The proposed framework naturally brings advantages such as in-context learning abilities and reusable self-supervised pretrained AudioMAE and latent diffusion models. Experiments on the major benchmarks of text-to-audio, text-to-music, and text-to-speech demonstrate state-of-the-art or competitive performance against previous approaches. Our code, pretrained model, and demo are available at [this https URL](https://audioldm.github.io/audioldm2).*
This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi). The original codebase can be
found at [haoheliu/audioldm2](https://github.com/haoheliu/audioldm2).
This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi). The original codebase can be found at [haoheliu/audioldm2](https://github.com/haoheliu/audioldm2).
## Tips
### Choosing a checkpoint
AudioLDM2 comes in three variants. Two of these checkpoints are applicable to the general task of text-to-audio
generation. The third checkpoint is trained exclusively on text-to-music generation.
AudioLDM2 comes in three variants. Two of these checkpoints are applicable to the general task of text-to-audio generation. The third checkpoint is trained exclusively on text-to-music generation.
All checkpoints share the same model size for the text encoders and VAE. They differ in the size and depth of the UNet.
All checkpoints share the same model size for the text encoders and VAE. They differ in the size and depth of the UNet.
See table below for details on the three checkpoints:
| Checkpoint | Task | UNet Model Size | Total Model Size | Training Data / h |
@@ -54,7 +41,7 @@ See table below for details on the three checkpoints:
* Descriptive prompt inputs work best: use adjectives to describe the sound (e.g. "high quality" or "clear") and make the prompt context specific (e.g. "water stream in a forest" instead of "stream").
* It's best to use general terms like "cat" or "dog" instead of specific names or abstract objects the model may not be familiar with.
* Using a **negative prompt** can significantly improve the quality of the generated waveform, by guiding the generation away from terms that correspond to poor quality audio. Try using a negative prompt of "Low quality."
* Using a **negative prompt** can significantly improve the quality of the generated waveform, by guiding the generation away from terms that correspond to poor quality audio. Try using a negative prompt of "Low quality."
### Controlling inference
@@ -63,16 +50,14 @@ See table below for details on the three checkpoints:
### Evaluating generated waveforms:
* The quality of the generated waveforms can vary significantly based on the seed. Try generating with different seeds until you find a satisfactory generation
* The quality of the generated waveforms can vary significantly based on the seed. Try generating with different seeds until you find a satisfactory generation.
* Multiple waveforms can be generated in one go: set `num_waveforms_per_prompt` to a value greater than 1. Automatic scoring will be performed between the generated waveforms and prompt text, and the audios ranked from best to worst accordingly.
The following example demonstrates how to construct good music generation using the aforementioned tips: [example](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2Pipeline.__call__.example).
<Tip>
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between
scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines)
section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
@@ -90,4 +75,4 @@ section to learn how to efficiently load the same components into multiple pipel
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# BLIP-Diffusion
BLIP-Diffusion was proposed in [BLIP-Diffusion: Pre-trained Subject Representation for Controllable Text-to-Image Generation and Editing](https://arxiv.org/abs/2305.14720). It enables zero-shot subject-driven generation and control-guided zero-shot generation.
The abstract from the paper is:
*Subject-driven text-to-image generation models create novel renditions of an input subject based on text prompts. Existing models suffer from lengthy fine-tuning and difficulties preserving the subject fidelity. To overcome these limitations, we introduce BLIP-Diffusion, a new subject-driven image generation model that supports multimodal control which consumes inputs of subject images and text prompts. Unlike other subject-driven generation models, BLIP-Diffusion introduces a new multimodal encoder which is pre-trained to provide subject representation. We first pre-train the multimodal encoder following BLIP-2 to produce visual representation aligned with the text. Then we design a subject representation learning task which enables a diffusion model to leverage such visual representation and generates new subject renditions. Compared with previous methods such as DreamBooth, our model enables zero-shot subject-driven generation, and efficient fine-tuning for customized subject with up to 20x speedup. We also demonstrate that BLIP-Diffusion can be flexibly combined with existing techniques such as ControlNet and prompt-to-prompt to enable novel subject-driven generation and editing applications. Project page at [this https URL](https://dxli94.github.io/BLIP-Diffusion-website/).*
The original codebase can be found at [salesforce/LAVIS](https://github.com/salesforce/LAVIS/tree/main/projects/blip-diffusion). You can find the official BLIP-Diffusion checkpoints under the [hf.co/SalesForce](https://hf.co/SalesForce) organization.
`BlipDiffusionPipeline` and `BlipDiffusionControlNetPipeline` were contributed by [`ayushtues`](https://github.com/ayushtues/).
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Consistency Models
Consistency Models were proposed in [Consistency Models](https://huggingface.co/papers/2303.01469) by Yang Song, Prafulla Dhariwal, Mark Chen, and Ilya Sutskever.
The abstract from the paper is:
*Diffusion models have significantly advanced the fields of image, audio, and video generation, but they depend on an iterative sampling process that causes slow generation. To overcome this limitation, we propose consistency models, a new family of models that generate high quality samples by directly mapping noise to data. They support fast one-step generation by design, while still allowing multistep sampling to trade compute for sample quality. They also support zero-shot data editing, such as image inpainting, colorization, and super-resolution, without requiring explicit training on these tasks. Consistency models can be trained either by distilling pre-trained diffusion models, or as standalone generative models altogether. Through extensive experiments, we demonstrate that they outperform existing distillation techniques for diffusion models in one- and few-step sampling, achieving the new state-of-the-art FID of 3.55 on CIFAR-10 and 6.20 on ImageNet 64x64 for one-step generation. When trained in isolation, consistency models become a new family of generative models that can outperform existing one-step, non-adversarial generative models on standard benchmarks such as CIFAR-10, ImageNet 64x64 and LSUN 256x256.*
*Diffusion models have significantly advanced the fields of image, audio, and video generation, but they depend on an iterative sampling process that causes slow generation. To overcome this limitation, we propose consistency models, a new family of models that generate high quality samples by directly mapping noise to data. They support fast one-step generation by design, while still allowing multistep sampling to trade compute for sample quality. They also support zero-shot data editing, such as image inpainting, colorization, and super-resolution, without requiring explicit training on these tasks. Consistency models can be trained either by distilling pre-trained diffusion models, or as standalone generative models altogether. Through extensive experiments, we demonstrate that they outperform existing distillation techniques for diffusion models in one- and few-step sampling, achieving the new state-of-the-art FID of 3.55 on CIFAR-10 and 6.20 on ImageNet 64x64 for one-step generation. When trained in isolation, consistency models become a new family of generative models that can outperform existing one-step, non-adversarial generative models on standard benchmarks such as CIFAR-10, ImageNet 64x64 and LSUN 256x256.*
The original codebase can be found at [openai/consistency_models](https://github.com/openai/consistency_models), and additional checkpoints are available at [openai](https://huggingface.co/openai).
@@ -27,17 +39,18 @@ For an additional speed-up, use `torch.compile` to generate multiple images in <
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,13 +12,13 @@ specific language governing permissions and limitations under the License.
# ControlNet
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang and Maneesh Agrawala.
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
The abstract from the paper is:
*We present a neural network structure, ControlNet, to control pretrained large diffusion models to support additional input conditions. The ControlNet learns task-specific conditions in an end-to-end way, and the learning is robust even when the training dataset is small (< 50k). Moreover, training a ControlNet is as fast as fine-tuning a diffusion model, and the model can be trained on a personal devices. Alternatively, if powerful computation clusters are available, the model can scale to large amounts (millions to billions) of data. We report that large diffusion models like Stable Diffusion can be augmented with ControlNets to enable conditional inputs like edge maps, segmentation maps, keypoints, etc. This may enrich the methods to control large diffusion models and further facilitate related applications.*
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
This model was contributed by [takuma104](https://huggingface.co/takuma104). ❤️
@@ -26,7 +26,7 @@ The original codebase can be found at [lllyasviel/ControlNet](https://github.com
<Tip>
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
@@ -67,7 +67,6 @@ Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to le
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,13 +12,13 @@ specific language governing permissions and limitations under the License.
# ControlNet with Stable Diffusion XL
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang and Maneesh Agrawala.
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
The abstract from the paper is:
*We present a neural network structure, ControlNet, to control pretrained large diffusion models to support additional input conditions. The ControlNet learns task-specific conditions in an end-to-end way, and the learning is robust even when the training dataset is small (< 50k). Moreover, training a ControlNet is as fast as fine-tuning a diffusion model, and the model can be trained on a personal devices. Alternatively, if powerful computation clusters are available, the model can scale to large amounts (millions to billions) of data. We report that large diffusion models like Stable Diffusion can be augmented with ControlNets to enable conditional inputs like edge maps, segmentation maps, keypoints, etc. This may enrich the methods to control large diffusion models and further facilitate related applications.*
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
You can find additional smaller Stable Diffusion XL (SDXL) ControlNet checkpoints from the 🤗 [Diffusers](https://huggingface.co/diffusers) Hub organization, and browse [community-trained](https://huggingface.co/models?other=stable-diffusion-xl&other=controlnet) checkpoints on the Hub.
@@ -28,11 +28,11 @@ You can find additional smaller Stable Diffusion XL (SDXL) ControlNet checkpoint
</Tip>
If you don't see a checkpoint you're interested in, you can train your own SDXL ControlNet with our [training script](https://github.com/huggingface/diffusers/blob/main/examples/controlnet/README_sdxl.md).
If you don't see a checkpoint you're interested in, you can train your own SDXL ControlNet with our [training script](../../../../../examples/controlnet/README_sdxl).
<Tip>
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
@@ -41,6 +41,15 @@ Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to le
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Cycle Diffusion
Cycle Diffusion is a text guided image-to-image generation model proposed in [Unifying Diffusion Models' Latent Space, with Applications to CycleDiffusion and Guidance](https://huggingface.co/papers/2210.05559) by Chen Henry Wu, Fernando De la Torre.
The abstract from the paper is:
*Diffusion models have achieved unprecedented performance in generative modeling. The commonly-adopted formulation of the latent code of diffusion models is a sequence of gradually denoised samples, as opposed to the simpler (e.g., Gaussian) latent space of GANs, VAEs, and normalizing flows. This paper provides an alternative, Gaussian formulation of the latent space of various diffusion models, as well as an invertible DPM-Encoder that maps images into the latent space. While our formulation is purely based on the definition of diffusion models, we demonstrate several intriguing consequences. (1) Empirically, we observe that a common latent space emerges from two diffusion models trained independently on related domains. In light of this finding, we propose CycleDiffusion, which uses DPM-Encoder for unpaired image-to-image translation. Furthermore, applying CycleDiffusion to text-to-image diffusion models, we show that large-scale text-to-image diffusion models can be used as zero-shot image-to-image editors. (2) One can guide pre-trained diffusion models and GANs by controlling the latent codes in a unified, plug-and-play formulation based on energy-based models. Using the CLIP model and a face recognition model as guidance, we demonstrate that diffusion models have better coverage of low-density sub-populations and individuals than GANs.*
<Tip>
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -16,11 +16,10 @@ specific language governing permissions and limitations under the License.
Dance Diffusion is the first in a suite of generative audio tools for producers and musicians released by [Harmonai](https://github.com/Harmonai-org).
The original codebase of this implementation can be found at [Harmonai-org](https://github.com/Harmonai-org/sample-generator).
<Tip>
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
@@ -30,4 +29,4 @@ Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to le
- __call__
## AudioPipelineOutput
[[autodoc]] pipelines.AudioPipelineOutput
[[autodoc]] pipelines.AudioPipelineOutput
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