* find & replace all FloatTensors to Tensor
* apply formatting
* Update torch.FloatTensor to torch.Tensor in the remaining files
* formatting
* Fix the rest of the places where FloatTensor is used as well as in documentation
* formatting
* Update new file from FloatTensor to Tensor
* Remove dead code
* PylancereportGeneralTypeIssues: Strings nested within an f-string cannot use the same quote character as the f-string prior to Python 3.12.
* Remove dead code
SDXL LoRA weights for text encoders should be decoupled on save
The method checks if at least one of unet, text_encoder and
text_encoder_2 lora weights are passed, which was not reflected in the
implentation.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
`model_output.shape` may only have rank 1.
There are warnings related to use of random keys.
```
tests/schedulers/test_scheduler_flax.py: 13 warnings
/Users/phillypham/diffusers/src/diffusers/schedulers/scheduling_ddpm_flax.py:268: FutureWarning: normal accepts a single key, but was given a key array of shape (1, 2) != (). Use jax.vmap for batching. In a future JAX version, this will be an error.
noise = jax.random.normal(split_key, shape=model_output.shape, dtype=self.dtype)
tests/schedulers/test_scheduler_flax.py::FlaxDDPMSchedulerTest::test_betas
/Users/phillypham/virtualenv/diffusers/lib/python3.9/site-packages/jax/_src/random.py:731: FutureWarning: uniform accepts a single key, but was given a key array of shape (1,) != (). Use jax.vmap for batching. In a future JAX version, this will be an error.
u = uniform(key, shape, dtype, lo, hi) # type: ignore[arg-type]
```
* 7879 - adjust documentation to use naruto dataset, since pokemon is now gated
* replace references to pokemon in docs
* more references to pokemon replaced
* Japanese translation update
---------
Co-authored-by: bghira <bghira@users.github.com>
* Add Ascend NPU support for SDXL fine-tuning and fix the model saving bug when using DeepSpeed.
* fix check code quality
* Decouple the NPU flash attention and make it an independent module.
* add doc and unit tests for npu flash attention.
---------
Co-authored-by: mhh001 <mahonghao1@huawei.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* chore: reducing model sizes
* chore: shrinks further
* chore: shrinks further
* chore: shrinking model for img2img pipeline
* chore: reducing size of model for inpaint pipeline
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
FlaxStableDiffusionSafetyChecker sets main_input_name to "clip_input".
This makes StableDiffusionSafetyChecker consistent.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Added get_velocity function to EulerDiscreteScheduler.
* Fix white space on blank lines
* Added copied from statement
* back to the original.
---------
Co-authored-by: Ruining Li <ruining@robots.ox.ac.uk>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
swap the order for do_classifier_free_guidance concat with repeat
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
* Check for latents, before calling prepare_latents - sdxlImg2Img
* Added latents check for all the img2img pipeline
* Fixed silly mistake while checking latents as None
A new function compute_dream_and_update_latents has been added to the
training utilities that allows you to do DREAM rectified training in line
with the paper https://arxiv.org/abs/2312.00210. The method can be used
with an extra argument in the train_text_to_image.py script.
Co-authored-by: Jimmy <39@🇺🇸.com>
* Convert channel order to BGR for the watermark encoder. Convert the watermarked BGR images back to RGB. Fixes#6292
* Revert channel order before stacking images to overcome limitations that negative strides are currently not supported
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Fixed wrong decorator by modifying it to @classmethod.
* Updated the method and it's argument.
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* add scheduled pseudo-huber loss training scripts
See #7488
* add reduction modes to huber loss
* [DB Lora] *2 multiplier to huber loss cause of 1/2 a^2 conv.
pairing of c6495def1f
* [DB Lora] add option for smooth l1 (huber / delta)
Pairing of dd22958caa
* [DB Lora] unify huber scheduling
Pairing of 19a834c3ab
* [DB Lora] add snr huber scheduler
Pairing of 47fb1a6854
* fixup examples link
* use snr schedule by default in DB
* update all huber scripts with snr
* code quality
* huber: make style && make quality
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Initialize target_unet from unet rather than teacher_unet so that we correctly add time_embedding.cond_proj if necessary.
* Use UNet2DConditionModel.from_config to initialize target_unet from unet's config.
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* give it a shot.
* print.
* correct assertion.
* gather results from the rest of the tests.
* change the assertion values where needed.
* remove print statements.
* get device <-> component mapping when using multiple gpus.
* condition the device_map bits.
* relax condition
* device_map progress.
* device_map enhancement
* some cleaning up and debugging
* Apply suggestions from code review
Co-authored-by: Marc Sun <57196510+SunMarc@users.noreply.github.com>
* incorporate suggestions from PR.
* remove multi-gpu condition for now.
* guard check the component -> device mapping
* fix: device_memory variable
* dispatching transformers model to have force_hooks=True
* better guarding for transformers device_map
* introduce support balanced_low_memory and balanced_ultra_low_memory.
* remove device_map patch.
* fix: intermediate variable scoping.
* fix: condition in cpu offload.
* fix: flax class restrictions.
* remove modifications from cpu_offload and model_offload
* incorporate changes.
* add a simple forward pass test
* add: torch_device in get_inputs()
* add: tests
* remove print
* safe-guard to(), model offloading and cpu offloading when balanced is used as a device_map.
* style
* remove .
* safeguard device_map with more checks and remove invalid device_mapping strategues.
* make a class attribute and adjust tests accordingly.
* fix device_map check
* fix test
* adjust comment
* fix: device_map attribute
* fix: dispatching.
* max_memory test for pipeline
* version guard the tests
* fix guard.
* address review feedback.
* reset_device_map method.
* add: test for reset_hf_device_map
* fix a couple things.
* add reset_device_map() in the error message.
* add tests for checking reset_device_map doesn't have unintended consequences.
* fix reset_device_map and offloading tests.
* create _get_final_device_map utility.
* hf_device_map -> _hf_device_map
* add documentation
* add notes suggested by Marc.
* styling.
* Apply suggestions from code review
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* move updates within gpu condition.
* other docs related things
* note on ignore a device not specified in .
* provide a suggestion if device mapping errors out.
* fix: typo.
* _hf_device_map -> hf_device_map
* Empty-Commit
* add: example hf_device_map.
---------
Co-authored-by: Marc Sun <57196510+SunMarc@users.noreply.github.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* remove libsndfile1-dev and libgl1 from workflows and ensure that re present in the respective dockerfiles.
* change to self-hosted runner; let's see 🤞
* add libsndfile1-dev libgl1 for now
* use self-hosted runners for building and push too.
* Restore unet params back to normal from EMA when validation call is finished
* empty commit
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Allow safety and feature extractor arguments to be passed to convert_from_ckpt
Allows management of safety checker and feature extractor
from outside of the convert ckpt class.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* reduce block sizes for unet1d.
* reduce blocks for unet_2d.
* reduce block size for unet_motion
* increase channels.
* correctly increase channels.
* reduce number of layers in unet2dconditionmodel tests.
* reduce block sizes for unet2dconditionmodel tests
* reduce block sizes for unet3dconditionmodel.
* fix: test_feed_forward_chunking
* fix: test_forward_with_norm_groups
* skip spatiotemporal tests on MPS.
* reduce block size in AutoencoderKL.
* reduce block sizes for vqmodel.
* further reduce block size.
* make style.
* Empty-Commit
* reduce sizes for ConsistencyDecoderVAETests
* further reduction.
* further block reductions in AutoencoderKL and AssymetricAutoencoderKL.
* massively reduce the block size in unet2dcontionmodel.
* reduce sizes for unet3d
* fix tests in unet3d.
* reduce blocks further in motion unet.
* fix: output shape
* add attention_head_dim to the test configuration.
* remove unexpected keyword arg
* up a bit.
* groups.
* up again
* fix
* Skip `test_freeu_enabled ` on MPS
* Small fixes
- import skip_mps correctly
- disable all instances of test_freeu_enabled
* Empty commit to trigger tests
* Empty commit to trigger CI
* increase number of workers for the tests.
* move to beefier runner.
* improve the fast push tests too.
* use a beefy machine for pytorch pipeline tests
* up the number of workers further.
* UniPC Multistep add `rescale_betas_zero_snr`
Same patch as DPM and Euler with the patched final alpha cumprod
BF16 doesn't seem to break down, I think cause UniPC upcasts during some
phases already? We could still force an upcast since it only
loses ≈ 0.005 it/s for me but the difference in output is very small. A
better endeavor might upcasting in step() and removing all the other
upcasts elsewhere?
* UniPC ZSNR UT
* Re-add `rescale_betas_zsnr` doc oops
* UniPC UTs iterate solvers on FP16
It wasn't catching errs on order==3. Might be excessive?
* UniPC Multistep fix tensor dtype/device on order=3
* UniPC UTs Add v_pred to fp16 test iter
For completions sake. Probably overkill?
* 7529 do not disable autocast for cuda devices
* Remove typecasting error check for non-mps platforms, as a correct autocast implementation makes it a non-issue
* add autocast fix to other training examples
* disable native_amp for dreambooth (sdxl)
* disable native_amp for pix2pix (sdxl)
* remove tests from remaining files
* disable native_amp on huggingface accelerator for every training example that uses it
* convert more usages of autocast to nullcontext, make style fixes
* make style fixes
* style.
* Empty-Commit
---------
Co-authored-by: bghira <bghira@users.github.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* start printing the tensors.
* print full throttle
* set static slices for 7 tests.
* remove printing.
* flatten
* disable test for controlnet
* what happens when things are seeded properly?
* set the right value
* style./
* make pia test fail to check things
* print.
* fix pia.
* checking for animatediff.
* fix: animatediff.
* video synthesis
* final piece.
* style.
* print guess.
* fix: assertion for control guess.
---------
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
* Add `final_sigma_zero` to UniPCMultistep
Effectively the same trick as DDIM's `set_alpha_to_one` and
DPM's `final_sigma_type='zero'`.
Currently False by default but maybe this should be True?
* `final_sigma_zero: bool` -> `final_sigmas_type: str`
Should 1:1 match DPM Multistep now.
* Set `final_sigmas_type='sigma_min'` in UniPC UTs
* Initial commit
* Implemented block lora
- implemented block lora
- updated docs
- added tests
* Finishing up
* Reverted unrelated changes made by make style
* Fixed typo
* Fixed bug + Made text_encoder_2 scalable
* Integrated some review feedback
* Incorporated review feedback
* Fix tests
* Made every module configurable
* Adapter to new lora test structure
* Final cleanup
* Some more final fixes
- Included examples in `using_peft_for_inference.md`
- Added hint that only attns are scaled
- Removed NoneTypes
- Added test to check mismatching lens of adapter names / weights raise error
* Update using_peft_for_inference.md
* Update using_peft_for_inference.md
* Make style, quality, fix-copies
* Updated tutorial;Warning if scale/adapter mismatch
* floats are forwarded as-is; changed tutorial scale
* make style, quality, fix-copies
* Fixed typo in tutorial
* Moved some warnings into `lora_loader_utils.py`
* Moved scale/lora mismatch warnings back
* Integrated final review suggestions
* Empty commit to trigger CI
* Reverted emoty commit to trigger CI
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* speed up test_vae_slicing in animatediff
* speed up test_karras_schedulers_shape for attend and excite.
* style.
* get the static slices out.
* specify torch print options.
* modify
* test run with controlnet
* specify kwarg
* fix: things
* not None
* flatten
* controlnet img2img
* complete controlet sd
* finish more
* finish more
* finish more
* finish more
* finish the final batch
* add cpu check for expected_pipe_slice.
* finish the rest
* remove print
* style
* fix ssd1b controlnet test
* checking ssd1b
* disable the test.
* make the test_ip_adapter_single controlnet test more robust
* fix: simple inpaint
* multi
* disable panorama
* enable again
* panorama is shaky so leave it for now
* remove print
* raise tolerance.
* Bug fix for controlnetpipeline check_image
Bug fix for controlnetpipeline check_image when using multicontrolnet and prompt list
* Update test_inference_multiple_prompt_input function
* Update test_controlnet.py
add test for multiple prompts and multiple image conditioning
* Update test_controlnet.py
Fix format error
---------
Co-authored-by: Lvkesheng Shen <45848260+Fantast416@users.noreply.github.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* add remove_all_hooks
* a few more fix and tests
* up
* Update src/diffusers/pipelines/pipeline_utils.py
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* split tests
* add
---------
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* apple mps: training support for SDXL LoRA
* sdxl: support training lora, dreambooth, t2i, pix2pix, and controlnet on apple mps
---------
Co-authored-by: bghira <bghira@users.github.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* mps: fix XL pipeline inference at training time due to upstream pytorch bug
* Update src/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl.py
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* apply the safe-guarding logic elsewhere.
---------
Co-authored-by: bghira <bghira@users.github.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
you cannot specify `type="bool"` and `action="store_true"` at the same time.
remove excessive and buggy `type=bool`.
Co-authored-by: Linoy Tsaban <57615435+linoytsaban@users.noreply.github.com>
* feat: support dora loras from community
* safe-guard dora operations under peft version.
* pop use_dora when False
* make dora lora from kohya work.
* fix: kohya conversion utils.
* add a fast test for DoRA compatibility..
* add a nightly test.
* fixed typo
* updated doc to be consistent in naming
* make style/quality
* preprocessing for 4 channels and not 6
* make style
* test for 4c
* make style/quality
* fixed test on cpu
* fixed doc typo
* changed default ckpt to 4c
* Update pipeline_stable_diffusion_ldm3d.py
* fix bug
---------
Co-authored-by: Aflalo <estellea@isl-iam1.rr.intel.com>
Co-authored-by: Aflalo <estellea@isl-gpu33.rr.intel.com>
Co-authored-by: Aflalo <estellea@isl-gpu38.rr.intel.com>
* Add properties and `IPAdapterTesterMixin` tests for `StableDiffusionPanoramaPipeline`
* Update torch manual seed to use `torch.Generator(device=device)`
* Refactor 📞🔙 to support `callback_on_step_end`
* make fix-copies
* fix freeinit impl
* fix progress bar
* fix progress bar and remove old code
* fix num_inference_steps==1 case for freeinit by atleast running 1 step when fast sampling enabled
* checking to improve pipelines.
* more fixes.
* add: tip to encourage the usage of revision
* Apply suggestions from code review
* retrigger ci
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Fix ControlNetModel.from_unet do not load add_embedding
* delete white space in blank line
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* debugging
* let's see the numbers
* let's see the numbers
* let's see the numbers
* restrict tolerance.
* increase inference steps.
* shallow copy of cross_attentionkwargs
* remove print
* pop scale from the top-level unet instead of getting it.
* improve readability.
* Apply suggestions from code review
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* fix a little bit.
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Add properties and `IPAdapterTesterMixin` tests for `StableDiffusionPanoramaPipeline`
* Fix variable name typo and update comments
* Update deprecated `output_type="numpy"` to "np" in test files
* Discard changes to src/diffusers/pipelines/stable_diffusion_panorama/pipeline_stable_diffusion_panorama.py
* Update test_stable_diffusion_panorama.py
* Update numbers in README.md
* Update get_guidance_scale_embedding method to use timesteps instead of w
* Update number of checkpoints in README.md
* Add type hints and fix var name
* Fix PyTorch's convention for inplace functions
* Fix a typo
* Revert "Fix PyTorch's convention for inplace functions"
This reverts commit 74350cf65b.
* Fix typos
* Indent
* Refactor get_guidance_scale_embedding method in LEditsPPPipelineStableDiffusionXL class
* log loss per image
* add commandline param for per image loss logging
* style
* debug-loss -> debug_loss
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Change step_offset scheduler docstrings
* Mention it may be needed by some models
* More docstrings
These ones failed literal S&R because I performed it case-sensitive
which is fun.
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* add: support for notifying maintainers about the nightly test status
* add: a tempoerary workflow for validation.
* cancel in progress.
* runs-on
* clean up
* add: peft dep
* change device.
* multiple edits.
* remove temp workflow.
* add: a workflow to check if docker containers can be built if the files are modified.
* type
* unify docker image build test and push
* make it run on prs too.
* check
* check
* check
* check again.
* remove docker test build file.
* remove extra dependencies./
* check
* Initial commit
* Removed copy hints, as in original SDXLControlNetPipeline
Removed copy hints, as in original SDXLControlNetPipeline, as the `make fix-copies` seems to have issues with the @property decorator.
* Reverted changes to ControlNetXS
* Addendum to: Removed changes to ControlNetXS
* Added test+docs for mixture of denoiser
* Update docs/source/en/using-diffusers/controlnet.md
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update docs/source/en/using-diffusers/controlnet.md
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
* Update unet_2d_blocks.py
similar doc-string add to have in the original diffusion repository.
* Fix bug for mention in this issue section #6901
* Update src/diffusers/schedulers/scheduling_ddim_flax.py
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* Fix linter
* Restore empty line
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* copied from for t2i pipelines without ip adapter support.
* two more pipelines with proper copied from comments.
* revert to the original implementation
* throw error when patch inputs and layernorm are provided for transformers2d.
* add comment on supported norm_types in transformers2d
* more check
* fix: norm _type handling
* [bug] Fix float/int guidance scale not working in `StableVideoDiffusionPipeline`
* Add test to disable CFG on SVD
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* support and example launch for sdxl turbo
* White space fixes
* Trailing whitespace character
* ruff format
* fix guidance_scale and steps for turbo mode
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Radames Ajna <radamajna@gmail.com>
* update svd docs
* fix example doc string
* update return type hints/docs
* update type hints
* Fix typos in pipeline_stable_video_diffusion.py
* make style && make fix-copies
* Update src/diffusers/pipelines/stable_video_diffusion/pipeline_stable_video_diffusion.py
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Update src/diffusers/pipelines/stable_video_diffusion/pipeline_stable_video_diffusion.py
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* update based on suggestion
---------
Co-authored-by: M. Tolga Cangöz <mtcangoz@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Enable FakeTensorMode for EulerDiscreteScheduler scheduler
PyTorch's FakeTensorMode does not support `.numpy()` or `numpy.array()`
calls.
This PR replaces `sigmas` numpy tensor by a PyTorch tensor equivalent
Repro
```python
with torch._subclasses.FakeTensorMode() as fake_mode, ONNXTorchPatcher():
fake_model = DiffusionPipeline.from_pretrained(model_name, low_cpu_mem_usage=False)
```
that otherwise would fail with
`RuntimeError: .numpy() is not supported for tensor subclasses.`
* Address comments
* add tags for diffusers training
* add tags for diffusers training
* add tags for diffusers training
* add tags for diffusers training
* add tags for diffusers training
* add tags for diffusers training
* add dora tags for drambooth lora scripts
* style
* add is_dora arg
* style
* add dora training feature to sd 1.5 script
* added notes about DoRA training
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* initial
* check_inputs fix to the rest of pipelines
* add fix for no cfg too
* use of variable
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Add copyright notice to relevant files and fix typos
* Set `timestep_spacing` parameter of `StableDiffusionXLPipeline`'s scheduler to `'trailing'`.
* Update `StableDiffusionXLPipeline.from_single_file` by including EulerAncestralDiscreteScheduler with `timestep_spacing="trailing"` param.
* Update model loading method in SDXL Turbo documentation
* move model helper function in pipeline to EfficiencyMixin
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* DPMMultistep rescale_betas_zero_snr
* DPM upcast samples in step()
* DPM rescale_betas_zero_snr UT
* DPMSolverMulti move sample upcast after model convert
Avoids having to re-use the dtype.
* Add a newline for Ruff
* log_validation unification for controlnet.
* additional fixes.
* remove print.
* better reuse and loading
* make final inference run conditional.
* Update examples/controlnet/README_sdxl.md
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* resize the control image in the snippet.
---------
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
* Make LoRACompatibleConv padding_mode work.
* Format code style.
* add fast test
* Update src/diffusers/models/lora.py
Simplify the code by patrickvonplaten.
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* code refactor
* apply patrickvonplaten suggestion to simplify the code.
* rm test_lora_layers_old_backend.py and add test case in test_lora_layers_peft.py
* update test case.
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* modulize log validation
* run make style and refactor wanddb support
* remove redundant initialization
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* make checkpoint_merger pipeline pass the "variant" argument to from_pretrained()
* make style
---------
Co-authored-by: Lincoln Stein <lstein@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* add stable_diffusion_xl_ipex community pipeline
* make style for code quality check
* update docs as suggested
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* standardize model card
* fix tags
* correct import styling and update tags
* run make style and make quality
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* feat: allow low_cpu_mem_usage in ip adapter loading
* reduce the number of device placements.
* documentation.
* throw low_cpu_mem_usage warning only once from the main entry point.
* use load_model_into_meta in single file utils
* propagate to autoencoder and controlnet.
* correct class name access behaviour.
* remove torch_dtype from load_model_into_meta; seems unncessary
* remove incorrect kwarg
* style to avoid extra unnecessary line breaks
* fix: bias loading bug
* fixes for SDXL
* apply changes to the conversion script to match single_file_utils.py
* do transpose to match the single file loading logic.
Remove <cat-toy> validation prompt from textual_inversion_sdxl.py
The `<cat-toy>` validation prompt is a default choice for the example task in the README. But no other part of `textual_inversion_sdxl.py` references the cat toy and `textual_inversion.py` has a default validation prompt of `None` as well.
So bring `textual_inversion_sdxl.py` in line with `textual_inversion.py` and change default validation prompt to `None`
* attention_head_dim
* debug
* print more info
* correct num_attention_heads behaviour
* down_block_num_attention_heads -> num_attention_heads.
* correct the image link in doc.
* add: deprecation for num_attention_head
* fix: test argument to use attention_head_dim
* more fixes.
* quality
* address comments.
* remove depcrecation.
* add: support for passing ip adapter image embeddings
* debugging
* make feature_extractor unloading conditioned on safety_checker
* better condition
* type annotation
* index to look into value slices
* more debugging
* debugging
* serialize embeddings dict
* better conditioning
* remove unnecessary prints.
* Update src/diffusers/loaders/ip_adapter.py
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* make fix-copies and styling.
* styling and further copy fixing.
* fix: check_inputs call in controlnet sdxl img2img pipeline
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* feat: standarize model card creation for dreambooth training.
* correct 'inference
* remove comments.
* take component out of kwargs
* style
* add: card template to have a leaner description.
* widget support.
* propagate changes to train_dreambooth_lora
* propagate changes to custom diffusion
* make widget properly type-annotated
* fix: callback function name is incorrect
On this tutorial there is a function defined and then used inside `callback_on_step_end` argument, but the name was not correct (mismatch)
* fix: typo in num_timestep (correct is num_timesteps)
fixed property name
* remove _to_tensor
* remove _to_tensor definition
* remove _collapse_frames_into_batch
* remove lora for not bloating the code.
* remove sample_size.
* simplify code a bit more
* ensure timesteps are always in tensor.
* Fix `AutoencoderTiny` with `use_slicing`
When using slicing with AutoencoderTiny, the encoder mistakenly encodes the entire batch for every image in the batch.
* Fixed formatting issue
* add noise_offset param
* micro conditioning - wip
* image processing adjusted and moved to support micro conditioning
* change time ids to be computed inside train loop
* change time ids to be computed inside train loop
* change time ids to be computed inside train loop
* time ids shape fix
* move token replacement of validation prompt to the same section of instance prompt and class prompt
* add offset noise to sd15 advanced script
* fix token loading during validation
* fix token loading during validation in sdxl script
* a little clean
* style
* a little clean
* style
* sdxl script - a little clean + minor path fix
sd 1.5 script - change default resolution value
* ad 1.5 script - minor path fix
* fix missing comma in code example in model card
* clean up commented lines
* style
* remove time ids computed outside training loop - no longer used now that we utilize micro-conditioning, as all time ids are now computed inside the training loop
* style
* [WIP] - added draft readme, building off of examples/dreambooth/README.md
* readme
* readme
* readme
* readme
* readme
* readme
* readme
* readme
* removed --crops_coords_top_left from CLI args
* style
* fix missing shape bug due to missing RGB if statement
* add blog mention at the start of the reamde as well
* Update examples/advanced_diffusion_training/README.md
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* change note to render nicely as well
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Fix bug in ResnetBlock2D.forward when not USE_PEFT_BACKEND and using scale_shift for time emb where the lora scale gets overwritten.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* fix minsnr implementation for v-prediction case
* format code
* always compute snr when snr_gamma is specified
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* feat: explicitly tag to diffusers when using push_to_hub
* remove tags.
* reset repo.
* Apply suggestions from code review
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix: tests
* fix: push_to_hub behaviour for tagging from save_pretrained
* Apply suggestions from code review
Co-authored-by: Lucain <lucainp@gmail.com>
* Apply suggestions from code review
Co-authored-by: Lucain <lucainp@gmail.com>
* import fixes.
* add library name to existing model card.
* add: standalone test for generate_model_card
* fix tests for standalone method
* moved library_name to a better place.
* merge create_model_card and generate_model_card.
* fix test
* address lucain's comments
* fix return identation
* Apply suggestions from code review
Co-authored-by: Lucain <lucainp@gmail.com>
* address further comments.
* Update src/diffusers/pipelines/pipeline_utils.py
Co-authored-by: Lucain <lucainp@gmail.com>
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Lucain <lucainp@gmail.com>
* initial commit for unconditional/class-conditional consistency training script
* make style
* Add entry for consistency training script in community README.
* Move consistency training script from community to research_projects/consistency_training
* Add requirements.txt and README to research_projects/consistency_training directory.
* Manually revert community README changes for consistency training.
* Fix path to script after moving script to research projects.
* Add option to load U-Net weights from pretrained model.
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* begin animatediff img2video and video2video
* revert animatediff to original implementation
* add img2video as pipeline
* update
* add vid2vid pipeline
* update imports
* update
* remove copied from line for check_inputs
* update
* update examples
* add multi-batch support
* fix __init__.py files
* move img2vid to community
* update community readme and examples
* fix
* make fix-copies
* add vid2vid batch params
* apply suggestions from review
Co-Authored-By: Dhruv Nair <dhruv.nair@gmail.com>
* add test for animatediff vid2vid
* torch.stack -> torch.cat
Co-Authored-By: Dhruv Nair <dhruv.nair@gmail.com>
* make style
* docs for vid2vid
* update
* fix prepare_latents
* fix docs
* remove img2vid
* update README to :main
* remove slow test
* refactor pipeline output
* update docs
* update docs
* merge community readme from :main
* final fix i promise
* add support for url in animatediff example
* update example
* update callbacks to latest implementation
* Update src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Update src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix merge
* Apply suggestions from code review
* remove callback and callback_steps as suggested in review
* Update tests/pipelines/animatediff/test_animatediff_video2video.py
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix import error caused due to unet refactor in #6630
* fix numpy import error after tensor2vid refactor in #6626
* make fix-copies
* fix numpy error
* fix progress bar test
---------
Co-authored-by: Dhruv Nair <dhruv.nair@gmail.com>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* sd1.5 support in separate script
A quick adaptation to support people interested in using this method on 1.5 models.
* sd15 prompt text encoding and unet conversions
as per @linoytsaban 's recommendations. Testing would be appreciated,
* Readability and quality improvements
Removed some mentions of SDXL, and some arguments that don't apply to sd 1.5, and cleaned up some comments.
* make style/quality commands
* tracker rename and run-it doc
* Update examples/advanced_diffusion_training/train_dreambooth_lora_sd15_advanced.py
* Update examples/advanced_diffusion_training/train_dreambooth_lora_sd15_advanced.py
---------
Co-authored-by: Linoy Tsaban <57615435+linoytsaban@users.noreply.github.com>
* move unets to module 🦋
* parameterize unet-level import.
* fix flax unet2dcondition model import
* models __init__
* mildly depcrecating models.unet_2d_blocks in favor of models.unets.unet_2d_blocks.
* noqa
* correct depcrecation behaviour
* inherit from the actual classes.
* Empty-Commit
* backwards compatibility for unet_2d.py
* backward compatibility for unet_2d_condition
* bc for unet_1d
* bc for unet_1d_blocks
* Fixed the bug related to saving DeepSpeed models.
* Add information about training SD models using DeepSpeed to the README.
* Apply suggestions from code review
---------
Co-authored-by: mhh001 <mahonghao1@huawei.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* - extract function for stage in UNet2DConditionModel init & forward
- Add new function get_mid_block() to unet_2d_blocks.py
* add type hint to get_mid_block aligned with get_up_block and get_down_block; rename _set_xxx function
* add type hint and use keyword arguments
* remove `copy from` in versatile diffusion
* add animatediff img2vid
* fix
* Update examples/community/README.md
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* fix code snippet between ip adapter face id and animatediff img2vid
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* [Fix] Multiple image conditionings in a single batch for `StableDiffusionControlNetPipeline`.
* Refactor `check_inputs` in `StableDiffusionControlNetPipeline` to avoid redundant codes.
* Make the behavior of MultiControlNetModel to be the same to the original ControlNetModel
* Keep the code change minimum for nested list support
* Add fast test `test_inference_nested_image_input`
* Remove redundant check for nested image condition in `check_inputs`
Remove `len(image) == len(prompt)` check out of `check_image()`
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Better `ValueError` message for incompatible nested image list size
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Fix syntax error in `check_inputs`
* Remove warning message for multi-ControlNets with multiple prompts
* Fix a typo in test_controlnet.py
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Add test case for multiple prompts, single image conditioning in `StableDiffusionMultiControlNetPipelineFastTests`
* Improved `ValueError` message for nested `controlnet_conditioning_scale`
* Documenting the behavior of image list as `StableDiffusionControlNetPipeline` input
---------
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* Fixes#6418 Advanced Dreambooth LoRa Training
* change order of import to fix nit
* fix nit, use cast_training_params
* remove torch.compile fix, will move to a new PR
* remove unnecessary import
* Enable image resizing to adjust its height and width in StableDiffusionXLInstructPix2PixPipeline
* Ensure that validation is performed at every 'validation_step', not at every step
* fix: training resume from fp16.
* add: comment
* remove residue from another branch.
* remove more residues.
* thanks to Younes; no hacks.
* style.
* clean things a bit and modularize _set_state_dict_into_text_encoder
* add comment about the fix detailed.
* support compile
* make style
* move unwrap_model inside function
* change unwrap call
* run make style
* Update examples/dreambooth/train_dreambooth.py
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Revert "Update examples/dreambooth/train_dreambooth.py"
This reverts commit 70ab09732e.
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Remove conversion to RGB
* Add a Conversion Function
* Add type hint for convert_method
* Update src/diffusers/utils/loading_utils.py
Update docstring
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
* Update docstring
* Optimize imports
* Optimize imports (2)
* Reformat code
---------
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* base template file - train_instruct_pix2pix.py
* additional import and parser argument requried for lora
* finetune only instructpix2pix model -- no need to include these layers
* inject lora layers
* freeze unet model -- only lora layers are trained
* training modifications to train only lora parameters
* store only lora parameters
* move train script to research project
* run quality and style code checks
* move train script to a new folder
* add README
* update README
* update references in README
---------
Co-authored-by: Rahul Raman <rahulraman@gmail.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* enable stable-xl textual inversion
* check if optimizer_2 exists
* check text_encoder_2 before using
* add textual inversion for sdxl in a single file
* fix style
* fix example style
* reset for error changes
* add readme for sdxl
* fix style
* disable autocast as it will cause cast error when weight_dtype=bf16
* fix spelling error
* fix style and readme and 8bit optimizer
* add README_sdxl.md link
* add tracker key on log_validation
* run style
* rm the second center crop
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* add tutorials to toctree.yml
* fix title
* fix words
* add overview ja
* fix diffusion to 拡散
* fix line 21
* add space
* delete supported pipline
* fix tutorial_overview.md
* fix space
* fix typo
* Delete docs/source/ja/tutorials/using_peft_for_inference.md
this file is not translated
* Delete docs/source/ja/tutorials/basic_training.md
this file is not translated
* Delete docs/source/ja/tutorials/autopipeline.md
this file is not translated
* fix toctree
* add: experimental script for diffusion dpo training.
* random_crop cli.
* fix: caption tokenization.
* fix: pixel_values index.
* fix: grad?
* debug
* fix: reduction.
* fixes in the loss calculation.
* style
* fix: unwrap call.
* fix: validation inference.
* add: initial sdxl script
* debug
* make sure images in the tuple are of same res
* fix model_max_length
* report print
* boom
* fix: numerical issues.
* fix: resolution
* comment about resize.
* change the order of the training transformation.
* save call.
* debug
* remove print
* manually detaching necessary?
* use the same vae for validation.
* add: readme.
* unwrap text encoder when saving hook only for full text encoder tuning
* unwrap text encoder when saving hook only for full text encoder tuning
* save embeddings in each checkpoint as well
* save embeddings in each checkpoint as well
* save embeddings in each checkpoint as well
* Update examples/advanced_diffusion_training/train_dreambooth_lora_sdxl_advanced.py
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* add documentation for DeepCache
* fix typo
* add wandb url for DeepCache
* fix some typos
* add item in _toctree.yml
* update formats for arguments
* Update deepcache.md
* Update docs/source/en/optimization/deepcache.md
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* add StableDiffusionXLPipeline in doc
* Separate SDPipeline and SDXLPipeline
* Add the paper link of ablation experiments for hyper-parameters
* Apply suggestions from code review
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: Steven Liu <59462357+stevhliu@users.noreply.github.com>
* Make WDS pipeline interpolation type configurable.
* Make the VAE encoding batch size configurable.
* Make lora_alpha and lora_dropout configurable for LCM LoRA scripts.
* Generalize scalings_for_boundary_conditions function and make the timestep scaling configurable.
* Make LoRA target modules configurable for LCM-LoRA scripts.
* Move resolve_interpolation_mode to src/diffusers/training_utils.py and make interpolation type configurable in non-WDS script.
* apply suggestions from review
* debug
* debug test_with_different_scales_fusion_equivalence
* use the right method.
* place it right.
* let's see.
* let's see again
* alright then.
* add a comment.
* I added a new doc string to the class. This is more flexible to understanding other developers what are doing and where it's using.
* Update src/diffusers/models/unet_2d_blocks.py
This changes suggest by maintener.
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update src/diffusers/models/unet_2d_blocks.py
Add suggested text
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Update unet_2d_blocks.py
I changed the Parameter to Args text.
* Update unet_2d_blocks.py
proper indentation set in this file.
* Update unet_2d_blocks.py
a little bit of change in the act_fun argument line.
* I run the black command to reformat style in the code
* Update unet_2d_blocks.py
similar doc-string add to have in the original diffusion repository.
* Batter way to write binarize function
* Solve check_code_quality error
* My mistake to run pull request but not reformated file
* Update image_processor.py
* remove extra variable and space
* Update image_processor.py
* Run ruff libarary to reformat my file
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
Co-authored-by: YiYi Xu <yixu310@gmail.com>
* add: test to check if peft loras are loadable in non-peft envs.
* add torch_device approrpiately.
* fix: get_dummy_inputs().
* test logits.
* rename
* debug
* debug
* fix: generator
* new assertion values after fixing the seed.
* shape
* remove print statements and settle this.
* to update values.
* change values when lora config is initialized under a fixed seed.
* update colab link
* update notebook link
* sanity restored by getting the exact same values without peft.
* change timesteps used to calculate snr when --with_prior_preservation is enabled
* change timesteps used to calculate snr when --with_prior_preservation is enabled (canonical script)
* style
* revert canonical script to before snr gamma change
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Add unload_ip_adapter method
* Update attn_processors with original layers
* Add test
* Use set_default_attn_processor
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
* Fix gradient-checkpointing option is ignored in SDXL+LoRA training. (#6388)
* Fix gradient-checkpointing option is ignored in SD+LoRA training.
* Fix gradient checkpoint is not applied to text encoders. (SDXL+LoRA)
---------
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
echo "Quality check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make style && make quality'" >> $GITHUB_STEP_SUMMARY
check_repository_consistency:
needs:check_code_quality
runs-on:ubuntu-latest
steps:
- uses:actions/checkout@v3
- name:Set up Python
uses:actions/setup-python@v4
with:
python-version:"3.8"
- name:Install dependencies
run:|
python -m pip install --upgrade pip
pip install .[quality]
- name:Check repo consistency
run:|
python utils/check_copies.py
python utils/check_dummies.py
make deps_table_check_updated
- name:Check if failure
if:${{ failure() }}
run:|
echo "Repo consistency check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make fix-copies'" >> $GITHUB_STEP_SUMMARY
echo "Quality check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make style && make quality'" >> $GITHUB_STEP_SUMMARY
check_repository_consistency:
needs:check_code_quality
runs-on:ubuntu-latest
steps:
- uses:actions/checkout@v3
- name:Set up Python
uses:actions/setup-python@v4
with:
python-version:"3.8"
- name:Install dependencies
run:|
python -m pip install --upgrade pip
pip install .[quality]
- name:Check repo consistency
run:|
python utils/check_copies.py
python utils/check_dummies.py
make deps_table_check_updated
- name:Check if failure
if:${{ failure() }}
run:|
echo "Repo consistency check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make fix-copies'" >> $GITHUB_STEP_SUMMARY
@@ -77,7 +77,7 @@ Please refer to the [How to use Stable Diffusion in Apple Silicon](https://huggi
## Quickstart
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 16000+ checkpoints):
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 22000+ checkpoints):
```python
fromdiffusersimportDiffusionPipeline
@@ -219,7 +219,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
- https://github.com/deep-floyd/IF
- https://github.com/bentoml/BentoML
- https://github.com/bmaltais/kohya_ss
- +7000 other amazing GitHub repositories 💪
- +9000 other amazing GitHub repositories 💪
Thank you for using us ❤️.
@@ -238,7 +238,7 @@ We also want to thank @heejkoo for the very helpful overview of papers, code and
```bibtex
@misc{von-platen-etal-2022-diffusers,
author={Patrick von Platen and Suraj Patil and Anton Lozhkov and Pedro Cuenca and Nathan Lambert and Kashif Rasul and Mishig Davaadorj and Thomas Wolf},
author={Patrick von Platen and Suraj Patil and Anton Lozhkov and Pedro Cuenca and Nathan Lambert and Kashif Rasul and Mishig Davaadorj and Dhruv Nair and Sayak Paul and William Berman and Yiyi Xu and Steven Liu and Thomas Wolf},
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,14 +12,18 @@ specific language governing permissions and limitations under the License.
# IP-Adapter
[IP-Adapter](https://hf.co/papers/2308.06721) is a lightweight adapter that enables prompting a diffusion model with an image. This method decouples the cross-attention layers of the image and text features. The image features are generated from an image encoder. Files generated from IP-Adapter are only ~100MBs.
[IP-Adapter](https://hf.co/papers/2308.06721) is a lightweight adapter that enables prompting a diffusion model with an image. This method decouples the cross-attention layers of the image and text features. The image features are generated from an image encoder.
<Tip>
Learn how to load an IP-Adapter checkpoint and image in the [IP-Adapter](../../using-diffusers/loading_adapters#ip-adapter) loading guide.
Learn how to load an IP-Adapter checkpoint and image in the IP-Adapter [loading](../../using-diffusers/loading_adapters#ip-adapter) guide, and you can see how to use it in the [usage](../../using-diffusers/ip_adapter) guide.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# PEFT
Diffusers supports loading adapters such as [LoRA](../../using-diffusers/loading_adapters) with the [PEFT](https://huggingface.co/docs/peft/index) library with the [`~loaders.peft.PeftAdapterMixin`] class. This allows modeling classes in Diffusers like [`UNet2DConditionModel`] to load an adapter.
<Tip>
Refer to the [Inference with PEFT](../../tutorials/using_peft_for_inference.md) tutorial for an overview of how to use PEFT in Diffusers for inference.
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -10,13 +10,124 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Single files
# Loading Pipelines and Models via `from_single_file`
Diffusers supports loading pretrained pipeline (or model) weights stored in a single file, such as a `ckpt` or `safetensors` file. These single file types are typically produced from community trained models. There are three classes for loading single file weights:
The `from_single_file` method allows you to load supported pipelines using a single checkpoint file as opposed to the folder format used by Diffusers. This is useful if you are working with many of the Stable Diffusion Web UI's (such as A1111) that extensively rely on a single file to distribute all the components of a diffusion model.
- [`FromSingleFileMixin`] supports loading pretrained pipeline weights stored in a single file, which can either be a `ckpt` or `safetensors` file.
- [`FromOriginalVAEMixin`] supports loading a pretrained [`AutoencoderKL`] from pretrained ControlNet weights stored in a single file, which can either be a `ckpt` or `safetensors` file.
- [`FromOriginalControlnetMixin`] supports loading pretrained ControlNet weights stored in a singlefile, which can either be a `ckpt` or `safetensors` file.
The `from_single_file` method also supports loading models in their originally distributed format. This means that supported models that have been finetuned with other services can be loaded directly into supported Diffusers model objects and pipelines.
## Pipelines that currently support `from_single_file` loading
- [`StableDiffusionPipeline`]
- [`StableDiffusionImg2ImgPipeline`]
- [`StableDiffusionInpaintPipeline`]
- [`StableDiffusionControlNetPipeline`]
- [`StableDiffusionControlNetImg2ImgPipeline`]
- [`StableDiffusionControlNetInpaintPipeline`]
- [`StableDiffusionUpscalePipeline`]
- [`StableDiffusionXLPipeline`]
- [`StableDiffusionXLImg2ImgPipeline`]
- [`StableDiffusionXLInpaintPipeline`]
- [`StableDiffusionXLInstructPix2PixPipeline`]
- [`StableDiffusionXLControlNetPipeline`]
- [`StableDiffusionXLKDiffusionPipeline`]
- [`LatentConsistencyModelPipeline`]
- [`LatentConsistencyModelImg2ImgPipeline`]
- [`StableDiffusionControlNetXSPipeline`]
- [`StableDiffusionXLControlNetXSPipeline`]
- [`LEditsPPPipelineStableDiffusion`]
- [`LEditsPPPipelineStableDiffusionXL`]
- [`PIAPipeline`]
## Models that currently support `from_single_file` loading
## Setting components in a Pipeline using `from_single_file`
Swap components of the pipeline by passing them directly to the `from_single_file` method. e.g If you would like use a different scheduler than the pipeline default.
## Using a Diffusers model repository to configure single file loading
Under the hood, `from_single_file` will try to determine a model repository to use to configure the components of the pipeline. You can also pass in a repository id to the `config` argument of the `from_single_file` method to explicitly set the repository to use.
## Override configuration options when using single file loading
Override the default model or pipeline configuration options when using `from_single_file` by passing in the relevant arguments directly to the `from_single_file` method. Any argument that is supported by the model or pipeline class can be configured in this way:
In the example above, since we explicitly passed `repo_id="segmind/SSD-1B"`, it will use this [configuration file](https://huggingface.co/segmind/SSD-1B/blob/main/unet/config.json) from the "unet" subfolder in `"segmind/SSD-1B"` to configure the unet component included in the checkpoint; Similarly, it will use the `config.json` file from `"vae"` subfolder to configure the vae model, `config.json` file from text_encoder folder to configure text_encoder and so on.
Note that most of the time you do not need to explicitly a `config` argument, `from_single_file` will automatically map the checkpoint to a repo id (we will discuss this in more details in next section). However, this can be useful in cases where model components might have been changed from what was originally distributed or in cases where a checkpoint file might not have the necessary metadata to correctly determine the configuration to use for the pipeline.
<Tip>
@@ -24,14 +135,114 @@ To learn more about how to load single file weights, see the [Load different Sta
</Tip>
## Working with local files
As of `diffusers>=0.28.0` the `from_single_file` method will attempt to configure a pipeline or model by first inferring the model type from the checkpoint file and then using the model type to determine the appropriate model repo configuration to use from the Hugging Face Hub. For example, any single file checkpoint based on the Stable Diffusion XL base model will use the [`stabilityai/stable-diffusion-xl-base-1.0`](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) model repo to configure the pipeline.
If you are working in an environment with restricted internet access, it is recommended to download the config files and checkpoints for the model to your preferred directory and pass the local paths to the `pretrained_model_link_or_path` and `config` arguments of the `from_single_file` method.
By default this will download the checkpoints and config files to the [Hugging Face Hub cache directory](https://huggingface.co/docs/huggingface_hub/en/guides/manage-cache). You can also specify a local directory to download the files to by passing the `local_dir` argument to the `hf_hub_download` and `snapshot_download` functions.
## Working with local files on file systems that do not support symlinking
By default the `from_single_file` method relies on the `huggingface_hub` caching mechanism to fetch and store checkpoints and config files for models and pipelines. If you are working with a file system that does not support symlinking, it is recommended that you first download the checkpoint file to a local directory and disable symlinking by passing the `local_dir_use_symlink=False` argument to the `hf_hub_download` and `snapshot_download` functions.
Disabling symlinking means that the `huggingface_hub` caching mechanism has no way to determine whether a file has already been downloaded to the local directory. This means that the `hf_hub_download` and `snapshot_download` functions will download files to the local directory each time they are executed. If you are disabling symlinking, it is recommended that you separate the model download and loading steps to avoid downloading the same file multiple times.
</Tip>
## Using the original configuration file of a model
If you would like to configure the parameters of the model components in the pipeline using the orignal YAML configuration file, you can pass a local path or url to the original configuration file to the `original_config` argument of the `from_single_file` method.
In the example above, the `original_config` file is only used to configure the parameters of the individual model components of the pipeline. For example it will be used to configure parameters such as the `in_channels` of the `vae` model and `unet` model. It is not used to determine the type of component objects in the pipeline.
<Tip>
When using `original_config` with local_files_only=True`, Diffusers will attempt to infer the components based on the type signatures of pipeline class, rather than attempting to fetch the pipeline config from the Hugging Face Hub. This is to prevent backwards breaking changes in existing code that might not be able to connect to the internet to fetch the necessary pipeline config files.
This is not as reliable as providing a path to a local config repo and might lead to errors when configuring the pipeline. To avoid this, please run the pipeline with `local_files_only=False` once to download the appropriate pipeline config files to the local cache.
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# Consistency Decoder
Consistency decoder can be used to decode the latents from the denoising UNet in the [`StableDiffusionPipeline`]. This decoder was introduced in the [DALL-E 3 technical report](https://openai.com/dall-e-3).
Consistency decoder can be used to decode the latents from the denoising UNet in the [`StableDiffusionPipeline`]. This decoder was introduced in the [DALL-E 3 technical report](https://openai.com/dall-e-3).
The original codebase can be found at [openai/consistencydecoder](https://github.com/openai/consistencydecoder).
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# UVit2DModel
The [U-ViT](https://hf.co/papers/2301.11093) model is a vision transformer (ViT) based UNet. This model incorporates elements from ViT (considers all inputs such as time, conditions and noisy image patches as tokens) and a UNet (long skip connections between the shallow and deep layers). The skip connection is important for predicting pixel-level features. An additional 3x3 convolutional block is applied prior to the final output to improve image quality.
The abstract from the paper is:
*Currently, applying diffusion models in pixel space of high resolution images is difficult. Instead, existing approaches focus on diffusion in lower dimensional spaces (latent diffusion), or have multiple super-resolution levels of generation referred to as cascades. The downside is that these approaches add additional complexity to the diffusion framework. This paper aims to improve denoising diffusion for high resolution images while keeping the model as simple as possible. The paper is centered around the research question: How can one train a standard denoising diffusion models on high resolution images, and still obtain performance comparable to these alternate approaches? The four main findings are: 1) the noise schedule should be adjusted for high resolution images, 2) It is sufficient to scale only a particular part of the architecture, 3) dropout should be added at specific locations in the architecture, and 4) downsampling is an effective strategy to avoid high resolution feature maps. Combining these simple yet effective techniques, we achieve state-of-the-art on image generation among diffusion models without sampling modifiers on ImageNet.*
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,10 +12,16 @@ specific language governing permissions and limitations under the License.
# aMUSEd
Amused is a lightweight text to image model based off of the [muse](https://arxiv.org/pdf/2301.00704.pdf) architecture. Amused is particularly useful in applications that require a lightweight and fast model such as generating many images quickly at once.
aMUSEd was introduced in [aMUSEd: An Open MUSE Reproduction](https://huggingface.co/papers/2401.01808) by Suraj Patil, William Berman, Robin Rombach, and Patrick von Platen.
Amused is a lightweight text to image model based off of the [MUSE](https://arxiv.org/abs/2301.00704) architecture. Amused is particularly useful in applications that require a lightweight and fast model such as generating many images quickly at once.
Amused is a vqvae token based transformer that can generate an image in fewer forward passes than many diffusion models. In contrast with muse, it uses the smaller text encoder CLIP-L/14 instead of t5-xxl. Due to its small parameter count and few forward pass generation process, amused can generate many images quickly. This benefit is seen particularly at larger batch sizes.
The abstract from the paper is:
*We present aMUSEd, an open-source, lightweight masked image model (MIM) for text-to-image generation based on MUSE. With 10 percent of MUSE's parameters, aMUSEd is focused on fast image generation. We believe MIM is under-explored compared to latent diffusion, the prevailing approach for text-to-image generation. Compared to latent diffusion, MIM requires fewer inference steps and is more interpretable. Additionally, MIM can be fine-tuned to learn additional styles with only a single image. We hope to encourage further exploration of MIM by demonstrating its effectiveness on large-scale text-to-image generation and releasing reproducible training code. We also release checkpoints for two models which directly produce images at 256x256 and 512x512 resolutions.*
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -25,6 +25,7 @@ The abstract of the paper is the following:
| Pipeline | Tasks | Demo
|---|---|:---:|
| [AnimateDiffPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff.py) | *Text-to-Video Generation with AnimateDiff* |
| [AnimateDiffVideoToVideoPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py) | *Video-to-Video Generation with AnimateDiff* |
## Available checkpoints
@@ -32,6 +33,8 @@ Motion Adapter checkpoints can be found under [guoyww](https://huggingface.co/gu
## Usage example
### AnimateDiffPipeline
AnimateDiff works with a MotionAdapter checkpoint and a Stable Diffusion model checkpoint. The MotionAdapter is a collection of Motion Modules that are responsible for adding coherent motion across image frames. These modules are applied after the Resnet and Attention blocks in Stable Diffusion UNet.
The following example demonstrates how to use a *MotionAdapter* checkpoint with Diffusers for inference based on StableDiffusion-1.4/1.5.
@@ -98,6 +101,161 @@ AnimateDiff tends to work better with finetuned Stable Diffusion models. If you
</Tip>
### AnimateDiffSDXLPipeline
AnimateDiff can also be used with SDXL models. This is currently an experimental feature as only a beta release of the motion adapter checkpoint is available.
prompt="a panda surfing in the ocean, realistic, high quality",
negative_prompt="low quality, worst quality",
num_inference_steps=20,
guidance_scale=8,
width=1024,
height=1024,
num_frames=16,
)
frames=output.frames[0]
export_to_gif(frames,"animation.gif")
```
### AnimateDiffVideoToVideoPipeline
AnimateDiff can also be used to generate visually similar videos or enable style/character/background or other edits starting from an initial video, allowing you to seamlessly explore creative possibilities.
alt="closeup of tony stark, robert downey jr, fireworks"
style="width: 300px;" />
</td>
</tr>
</table>
## Using Motion LoRAs
Motion LoRAs are a collection of LoRAs that work with the `guoyww/animatediff-motion-adapter-v1-5-2` checkpoint. These LoRAs are responsible for adding specific types of motion to the animations.
[FreeInit: Bridging Initialization Gap in Video Diffusion Models](https://arxiv.org/abs/2312.07537) by Tianxing Wu, Chenyang Si, Yuming Jiang, Ziqi Huang, Ziwei Liu.
FreeInit is an effective method that improves temporal consistency and overall quality of videos generated using video-diffusion-models without any addition training. It can be applied to AnimateDiff, ModelScope, VideoCrafter and various other video generation models seamlessly at inference time, and works by iteratively refining the latent-initialization noise. More details can be found it the paper.
The following example demonstrates the usage of FreeInit.
FreeInit is not really free - the improved quality comes at the cost of extra computation. It requires sampling a few extra times depending on the `num_iters` parameter that is set when enabling it. Setting the `use_fast_sampling` parameter to `True` can improve the overall performance (at the cost of lower quality compared to when `use_fast_sampling=False` but still better results than vanilla video generation models).
</Tip>
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
[AnimateLCM](https://animatelcm.github.io/) is a motion module checkpoint and an [LCM LoRA](https://huggingface.co/docs/diffusers/using-diffusers/inference_with_lcm_lora) that have been created using a consistency learning strategy that decouples the distillation of the image generation priors and the motion generation priors.
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -20,7 +20,8 @@ The abstract of the paper is the following:
*Although audio generation shares commonalities across different types of audio, such as speech, music, and sound effects, designing models for each type requires careful consideration of specific objectives and biases that can significantly differ from those of other types. To bring us closer to a unified perspective of audio generation, this paper proposes a framework that utilizes the same learning method for speech, music, and sound effect generation. Our framework introduces a general representation of audio, called "language of audio" (LOA). Any audio can be translated into LOA based on AudioMAE, a self-supervised pre-trained representation learning model. In the generation process, we translate any modalities into LOA by using a GPT-2 model, and we perform self-supervised audio generation learning with a latent diffusion model conditioned on LOA. The proposed framework naturally brings advantages such as in-context learning abilities and reusable self-supervised pretrained AudioMAE and latent diffusion models. Experiments on the major benchmarks of text-to-audio, text-to-music, and text-to-speech demonstrate state-of-the-art or competitive performance against previous approaches. Our code, pretrained model, and demo are available at [this https URL](https://audioldm.github.io/audioldm2).*
This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi). The original codebase can be found at [haoheliu/audioldm2](https://github.com/haoheliu/audioldm2).
This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi) and [Nguyễn Công Tú Anh](https://github.com/tuanh123789). The original codebase can be
found at [haoheliu/audioldm2](https://github.com/haoheliu/audioldm2).
## Tips
@@ -36,6 +37,8 @@ See table below for details on the three checkpoints:
@@ -53,7 +56,7 @@ See table below for details on the three checkpoints:
* The quality of the generated waveforms can vary significantly based on the seed. Try generating with different seeds until you find a satisfactory generation.
* Multiple waveforms can be generated in one go: set `num_waveforms_per_prompt` to a value greater than 1. Automatic scoring will be performed between the generated waveforms and prompt text, and the audios ranked from best to worst accordingly.
The following example demonstrates how to construct good music generation using the aforementioned tips: [example](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2Pipeline.__call__.example).
The following example demonstrates how to construct good music and speech generation using the aforementioned tips: [example](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2Pipeline.__call__.example).
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,42 +12,10 @@ specific language governing permissions and limitations under the License.
# AutoPipeline
`AutoPipeline` is designed to:
1. make it easy for you to load a checkpoint for a task without knowing the specific pipeline class to use
2. use multiple pipelines in your workflow
Based on the task, the `AutoPipeline` class automatically retrieves the relevant pipeline given the name or path to the pretrained weights with the `from_pretrained()` method.
To seamlessly switch between tasks with the same checkpoint without reallocating additional memory, use the `from_pipe()` method to transfer the components from the original pipeline to the new one.
The `AutoPipeline` is designed to make it easy to load a checkpoint for a task without needing to know the specific pipeline class. Based on the task, the `AutoPipeline` automatically retrieves the correct pipeline class from the checkpoint `model_index.json` file.
> [!TIP]
> Check out the [AutoPipeline](../../tutorials/autopipeline) tutorial to learn how to use this API!
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNet-XS
ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
@@ -12,5 +24,16 @@ Here's the overview from the [project page](https://vislearn.github.io/ControlNe
This model was contributed by [UmerHA](https://twitter.com/UmerHAdil). ❤️
<Tip>
> 🧠 Make sure to check out the Schedulers [guide](https://huggingface.co/docs/diffusers/main/en/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](https://huggingface.co/docs/diffusers/main/en/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNet-XS with Stable Diffusion XL
ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
@@ -12,4 +24,22 @@ Here's the overview from the [project page](https://vislearn.github.io/ControlNe
This model was contributed by [UmerHA](https://twitter.com/UmerHAdil). ❤️
> 🧠 Make sure to check out the Schedulers [guide](https://huggingface.co/docs/diffusers/main/en/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](https://huggingface.co/docs/diffusers/main/en/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
<Tip warning={true}>
🧪 Many of the SDXL ControlNet checkpoints are experimental, and there is a lot of room for improvement. Feel free to open an [Issue](https://github.com/huggingface/diffusers/issues/new/choose) and leave us feedback on how we can improve!
</Tip>
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# I2VGen-XL
[I2VGen-XL: High-Quality Image-to-Video Synthesis via Cascaded Diffusion Models](https://hf.co/papers/2311.04145.pdf) by Shiwei Zhang, Jiayu Wang, Yingya Zhang, Kang Zhao, Hangjie Yuan, Zhiwu Qin, Xiang Wang, Deli Zhao, and Jingren Zhou.
The abstract from the paper is:
*Video synthesis has recently made remarkable strides benefiting from the rapid development of diffusion models. However, it still encounters challenges in terms of semantic accuracy, clarity and spatio-temporal continuity. They primarily arise from the scarcity of well-aligned text-video data and the complex inherent structure of videos, making it difficult for the model to simultaneously ensure semantic and qualitative excellence. In this report, we propose a cascaded I2VGen-XL approach that enhances model performance by decoupling these two factors and ensures the alignment of the input data by utilizing static images as a form of crucial guidance. I2VGen-XL consists of two stages: i) the base stage guarantees coherent semantics and preserves content from input images by using two hierarchical encoders, and ii) the refinement stage enhances the video's details by incorporating an additional brief text and improves the resolution to 1280×720. To improve the diversity, we collect around 35 million single-shot text-video pairs and 6 billion text-image pairs to optimize the model. By this means, I2VGen-XL can simultaneously enhance the semantic accuracy, continuity of details and clarity of generated videos. Through extensive experiments, we have investigated the underlying principles of I2VGen-XL and compared it with current top methods, which can demonstrate its effectiveness on diverse data. The source code and models will be publicly available at [this https URL](https://i2vgen-xl.github.io/).*
The original codebase can be found [here](https://github.com/ali-vilab/i2vgen-xl/). The model checkpoints can be found [here](https://huggingface.co/ali-vilab/).
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines. Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section [here](../../using-diffusers/svd#reduce-memory-usage).
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the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
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# LEDITS++
LEDITS++ was proposed in [LEDITS++: Limitless Image Editing using Text-to-Image Models](https://huggingface.co/papers/2311.16711) by Manuel Brack, Felix Friedrich, Katharina Kornmeier, Linoy Tsaban, Patrick Schramowski, Kristian Kersting, Apolinário Passos.
The abstract from the paper is:
*Text-to-image diffusion models have recently received increasing interest for their astonishing ability to produce high-fidelity images from solely text inputs. Subsequent research efforts aim to exploit and apply their capabilities to real image editing. However, existing image-to-image methods are often inefficient, imprecise, and of limited versatility. They either require time-consuming fine-tuning, deviate unnecessarily strongly from the input image, and/or lack support for multiple, simultaneous edits. To address these issues, we introduce LEDITS++, an efficient yet versatile and precise textual image manipulation technique. LEDITS++'s novel inversion approach requires no tuning nor optimization and produces high-fidelity results with a few diffusion steps. Second, our methodology supports multiple simultaneous edits and is architecture-agnostic. Third, we use a novel implicit masking technique that limits changes to relevant image regions. We propose the novel TEdBench++ benchmark as part of our exhaustive evaluation. Our results demonstrate the capabilities of LEDITS++ and its improvements over previous methods. The project page is available at https://leditsplusplus-project.static.hf.space .*
<Tip>
You can find additional information about LEDITS++ on the [project page](https://leditsplusplus-project.static.hf.space/index.html) and try it out in a [demo](https://huggingface.co/spaces/editing-images/leditsplusplus).
</Tip>
<Tip warning={true}>
Due to some backward compatability issues with the current diffusers implementation of [`~schedulers.DPMSolverMultistepScheduler`] this implementation of LEdits++ can no longer guarantee perfect inversion.
This issue is unlikely to have any noticeable effects on applied use-cases. However, we provide an alternative implementation that guarantees perfect inversion in a dedicated [GitHub repo](https://github.com/ml-research/ledits_pp).
</Tip>
We provide two distinct pipelines based on different pre-trained models.
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