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12 Commits

Author SHA1 Message Date
Sayak Paul
965e52ce61 Release: v0.19.2 2023-07-28 23:33:51 +05:30
Sayak Paul
b1e52794a2 [Feat] Support SDXL Kohya-style LoRA (#4287)
* sdxl lora changes.

* better name replacement.

* better replacement.

* debugging

* debugging

* debugging

* debugging

* debugging

* remove print.

* print state dict keys.

* print

* distingisuih better

* debuggable.

* fxi: tyests

* fix: arg from training script.

* access from class.

* run style

* debug

* save intermediate

* some simplifications for SDXL LoRA

* styling

* unet config is not needed in diffusers format.

* fix: dynamic SGM block mapping for SDXL kohya loras (#4322)

* Use lora compatible layers for linear proj_in/proj_out (#4323)

* improve condition for using the sgm_diffusers mapping

* informative comment.

* load compatible keys and embedding layer maaping.

* Get SDXL 1.0 example lora to load

* simplify

* specif ranks and hidden sizes.

* better handling of k rank and hidden

* debug

* debug

* debug

* debug

* debug

* fix: alpha keys

* add check for handling LoRAAttnAddedKVProcessor

* sanity comment

* modifications for text encoder SDXL

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* denugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* up

* up

* up

* up

* up

* up

* unneeded comments.

* unneeded comments.

* kwargs for the other attention processors.

* kwargs for the other attention processors.

* debugging

* debugging

* debugging

* debugging

* improve

* debugging

* debugging

* more print

* Fix alphas

* debugging

* debugging

* debugging

* debugging

* debugging

* debugging

* clean up

* clean up.

* debugging

* fix: text

---------

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Batuhan Taskaya <batuhan@python.org>
2023-07-28 23:28:22 +05:30
Patrick von Platen
c3e3a1ee10 [SDXL] Make watermarker optional under certain circumstances to improve usability of SDXL 1.0 (#4346)
* improve sdxl

* more fixes

* improve sdxl

* improve sdxl

* improve sdxl

* finish
2023-07-28 23:28:11 +05:30
Patrick von Platen
9cde56a729 [ONNX] Don't download ONNX model by default (#4338)
* [Download] Don't download ONNX weights by default

* [Download] Don't download ONNX weights by default

* [Download] Don't download ONNX weights by default

* fix more

* finish

* finish

* finish
2023-07-28 23:27:58 +05:30
Patrick von Platen
c63d7cdba0 [SDXL Refiner] Fix refiner forward pass for batched input (#4327)
* fix_batch_xl

* Fix other pipelines as well

* up

* up

* Update tests/pipelines/stable_diffusion_xl/test_stable_diffusion_xl_inpaint.py

* sort

* up

* Finish it all up Co-authored-by: Bagheera <bghira@users.github.com>

* Co-authored-by: Bagheera bghira@users.github.com

* Co-authored-by: Bagheera <bghira@users.github.com>

* Finish it all up Co-authored-by: Bagheera <bghira@users.github.com>
2023-07-28 23:27:36 +05:30
Patrick von Platen
aa4634a7fa Release: v0.19.1 2023-07-27 20:00:43 +02:00
Patrick von Platen
0709650e9d [Local loading] Correct bug with local files only (#4318)
* [Local loading] Correct bug with local files only

* file not found error

* fix

* finish
2023-07-27 20:00:21 +02:00
YiYi Xu
a9829164f4 fix a bug in StableDiffusionUpscalePipeline when prompt is None (#4278)
* fix batch_size

* add test

---------

Co-authored-by: yiyixuxu <yixu310@gmail,com>
2023-07-27 20:00:12 +02:00
Duong A. Nguyen
49c95178ad Fix SDXL conversion from original to diffusers (#4280)
* fix sdxl conversion

* convention
2023-07-27 20:00:02 +02:00
Patrick von Platen
c2f755bc62 [Torch.compile] Fixes torch compile graph break (#4315)
* fix torch compile

* Fix all

* make style
2023-07-27 19:59:55 +02:00
YiYi Xu
2fb877b66c update Kandinsky doc (#4301)
* update doc

* fix an error in autopipe doc

---------

Co-authored-by: yiyixuxu <yixu310@gmail,com>
2023-07-27 19:59:50 +02:00
Patrick von Platen
ef9824f9f7 Release: v0.19.0 2023-07-26 21:03:45 +02:00
790 changed files with 17852 additions and 111442 deletions

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@@ -13,9 +13,8 @@ body:
*Give your issue a fitting title. Assume that someone which very limited knowledge of diffusers can understand your issue. Add links to the source code, documentation other issues, pull requests etc...*
- 2. If your issue is about something not working, **always** provide a reproducible code snippet. The reader should be able to reproduce your issue by **only copy-pasting your code snippet into a Python shell**.
*The community cannot solve your issue if it cannot reproduce it. If your bug is related to training, add your training script and make everything needed to train public. Otherwise, just add a simple Python code snippet.*
- 3. Add the **minimum** amount of code / context that is needed to understand, reproduce your issue.
- 3. Add the **minimum amount of code / context that is needed to understand, reproduce your issue**.
*Make the life of maintainers easy. `diffusers` is getting many issues every day. Make sure your issue is about one bug and one bug only. Make sure you add only the context, code needed to understand your issues - nothing more. Generally, every issue is a way of documenting this library, try to make it a good documentation entry.*
- 4. For issues related to community pipelines (i.e., the pipelines located in the `examples/community` folder), please tag the author of the pipeline in your issue thread as those pipelines are not maintained.
- type: markdown
attributes:
value: |
@@ -61,46 +60,21 @@ body:
All issues are read by one of the core maintainers, so if you don't know who to tag, just leave this blank and
a core maintainer will ping the right person.
Please tag a maximum of 2 people.
Please tag fewer than 3 people.
General library related questions: @patrickvonplaten and @sayakpaul
Questions on DiffusionPipeline (Saving, Loading, From pretrained, ...):
Questions on the training examples: @williamberman, @sayakpaul, @yiyixuxu
Questions on pipelines:
- Stable Diffusion @yiyixuxu @DN6 @patrickvonplaten @sayakpaul @patrickvonplaten
- Stable Diffusion XL @yiyixuxu @sayakpaul @DN6 @patrickvonplaten
- Kandinsky @yiyixuxu @patrickvonplaten
- ControlNet @sayakpaul @yiyixuxu @DN6 @patrickvonplaten
- T2I Adapter @sayakpaul @yiyixuxu @DN6 @patrickvonplaten
- IF @DN6 @patrickvonplaten
- Text-to-Video / Video-to-Video @DN6 @sayakpaul @patrickvonplaten
- Wuerstchen @DN6 @patrickvonplaten
- Other: @yiyixuxu @DN6
Questions on memory optimizations, LoRA, float16, etc.: @williamberman, @patrickvonplaten, and @sayakpaul
Questions on models:
- UNet @DN6 @yiyixuxu @sayakpaul @patrickvonplaten
- VAE @sayakpaul @DN6 @yiyixuxu @patrickvonplaten
- Transformers/Attention @DN6 @yiyixuxu @sayakpaul @DN6 @patrickvonplaten
Questions on schedulers: @patrickvonplaten and @williamberman
Questions on Schedulers: @yiyixuxu @patrickvonplaten
Questions on LoRA: @sayakpaul @patrickvonplaten
Questions on Textual Inversion: @sayakpaul @patrickvonplaten
Questions on Training:
- DreamBooth @sayakpaul @patrickvonplaten
- Text-to-Image Fine-tuning @sayakpaul @patrickvonplaten
- Textual Inversion @sayakpaul @patrickvonplaten
- ControlNet @sayakpaul @patrickvonplaten
Questions on Tests: @DN6 @sayakpaul @yiyixuxu
Questions on Documentation: @stevhliu
Questions on models and pipelines: @patrickvonplaten, @sayakpaul, and @williamberman
Questions on JAX- and MPS-related things: @pcuenca
Questions on audio pipelines: @DN6 @patrickvonplaten
Questions on audio pipelines: @patrickvonplaten, @kashif, and @sanchit-gandhi
Documentation: @stevhliu and @yiyixuxu
placeholder: "@Username ..."

View File

@@ -41,7 +41,7 @@ Core library:
- Schedulers: @williamberman and @patrickvonplaten
- Pipelines: @patrickvonplaten and @sayakpaul
- Training examples: @sayakpaul and @patrickvonplaten
- Docs: @stevhliu and @yiyixuxu
- Docs: @stevenliu and @yiyixu
- JAX and MPS: @pcuenca
- Audio: @sanchit-gandhi
- General functionalities: @patrickvonplaten and @sayakpaul

View File

@@ -26,8 +26,6 @@ jobs:
image-name:
- diffusers-pytorch-cpu
- diffusers-pytorch-cuda
- diffusers-pytorch-compile-cuda
- diffusers-pytorch-xformers-cuda
- diffusers-flax-cpu
- diffusers-flax-tpu
- diffusers-onnxruntime-cpu

View File

@@ -16,7 +16,7 @@ jobs:
install_libgl1: true
package: diffusers
notebook_folder: diffusers_doc
languages: en ko zh ja pt
languages: en ko zh
secrets:
token: ${{ secrets.HUGGINGFACE_PUSH }}

View File

@@ -15,4 +15,4 @@ jobs:
pr_number: ${{ github.event.number }}
install_libgl1: true
package: diffusers
languages: en ko zh ja pt
languages: en ko zh

View File

@@ -20,7 +20,7 @@ jobs:
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
python-version: "3.7"
- name: Install dependencies
run: |
python -m pip install --upgrade pip

View File

@@ -1,34 +0,0 @@
name: Run Flax dependency tests
on:
pull_request:
branches:
- main
push:
branches:
- main
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
jobs:
check_flax_dependencies:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install -e .
pip install "jax[cpu]>=0.2.16,!=0.3.2"
pip install "flax>=0.4.1"
pip install "jaxlib>=0.1.65"
pip install pytest
- name: Check for soft dependencies
run: |
pytest tests/others/test_dependencies.py

View File

@@ -20,7 +20,7 @@ jobs:
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
python-version: "3.7"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
@@ -38,7 +38,7 @@ jobs:
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
python-version: "3.7"
- name: Install dependencies
run: |
python -m pip install --upgrade pip

View File

@@ -1,67 +0,0 @@
name: Fast tests for PRs - PEFT backend
on:
pull_request:
branches:
- main
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
env:
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 4
MKL_NUM_THREADS: 4
PYTEST_TIMEOUT: 60
jobs:
run_fast_tests:
strategy:
fail-fast: false
matrix:
config:
- name: LoRA
framework: lora
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_lora
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
python -m pip install -U git+https://github.com/huggingface/transformers.git
python -m pip install -U git+https://github.com/huggingface/peft.git
- name: Environment
run: |
python utils/print_env.py
- name: Run fast PyTorch LoRA CPU tests with PEFT backend
if: ${{ matrix.config.framework == 'lora' }}
run: |
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v \
--make-reports=tests_${{ matrix.config.report }} \
tests/lora/test_lora_layers_peft.py

View File

@@ -34,11 +34,6 @@ jobs:
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_models_schedulers
- name: LoRA
framework: lora
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_lora
- name: Fast Flax CPU tests
framework: flax
runner: docker-cpu
@@ -72,7 +67,6 @@ jobs:
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install accelerate
- name: Environment
run: |
@@ -94,14 +88,6 @@ jobs:
--make-reports=tests_${{ matrix.config.report }} \
tests/models tests/schedulers tests/others
- name: Run fast PyTorch LoRA CPU tests
if: ${{ matrix.config.framework == 'lora' }}
run: |
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx and not Dependency" \
--make-reports=tests_${{ matrix.config.report }} \
tests/lora
- name: Run fast Flax TPU tests
if: ${{ matrix.config.framework == 'flax' }}
run: |
@@ -115,64 +101,7 @@ jobs:
run: |
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_${{ matrix.config.report }} \
examples/test_examples.py
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: pr_${{ matrix.config.report }}_test_reports
path: reports
run_staging_tests:
strategy:
fail-fast: false
matrix:
config:
- name: Hub tests for models, schedulers, and pipelines
framework: hub_tests_pytorch
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_hub
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
- name: Environment
run: |
python utils/print_env.py
- name: Run Hub tests for models, schedulers, and pipelines on a staging env
if: ${{ matrix.config.framework == 'hub_tests_pytorch' }}
run: |
HUGGINGFACE_CO_STAGING=true python -m pytest \
-m "is_staging_test" \
--make-reports=tests_${{ matrix.config.report }} \
tests
examples/test_examples.py
- name: Failure short reports
if: ${{ failure() }}

View File

@@ -1,32 +0,0 @@
name: Run Torch dependency tests
on:
pull_request:
branches:
- main
push:
branches:
- main
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
jobs:
check_torch_dependencies:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install -e .
pip install torch torchvision torchaudio
pip install pytest
- name: Check for soft dependencies
run: |
pytest tests/others/test_dependencies.py

View File

@@ -1,11 +1,10 @@
name: Slow Tests on main
name: Slow tests on main
on:
push:
branches:
- main
env:
DIFFUSERS_IS_CI: yes
HF_HOME: /mnt/cache
@@ -13,371 +12,104 @@ env:
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
PIPELINE_USAGE_CUTOFF: 50000
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cpu # this is a CPU image, but we need it to fetch the matrix
options: --shm-size "16gb" --ipc host
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
run: |
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
echo $matrix
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: test-pipelines.json
path: reports
torch_pipelines_cuda_tests:
name: Torch Pipelines CUDA Slow Tests
needs: setup_torch_cuda_pipeline_matrix
run_slow_tests:
strategy:
fail-fast: false
max-parallel: 1
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Slow PyTorch CUDA checkpoint tests on Ubuntu
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
tests/pipelines/${{ matrix.module }}
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
config:
- name: Slow PyTorch CUDA tests on Ubuntu
framework: pytorch
runner: docker-gpu
image: diffusers/diffusers-pytorch-cuda
report: torch_cuda
- name: Slow Flax TPU tests on Ubuntu
framework: flax
runner: docker-tpu
image: diffusers/diffusers-flax-tpu
report: flax_tpu
- name: Slow ONNXRuntime CUDA tests on Ubuntu
framework: onnxruntime
runner: docker-gpu
image: diffusers/diffusers-onnxruntime-cuda
report: onnx_cuda
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: pipeline_${{ matrix.module }}_test_reports
path: reports
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
torch_cuda_tests:
name: Torch CUDA Tests
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
image: ${{ matrix.config.image }}
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ ${{ matrix.config.runner == 'docker-tpu' && '--privileged' || '--gpus 0'}}
defaults:
run:
shell: bash
strategy:
matrix:
module: [models, schedulers, lora, others]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
if : ${{ matrix.config.runner == 'docker-gpu' }}
run: |
nvidia-smi
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Run slow PyTorch CUDA tests
if: ${{ matrix.config.framework == 'pytorch' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_cuda \
tests/${{ matrix.module }}
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_cuda_stats.txt
cat reports/tests_torch_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_cuda_test_reports
path: reports
peft_cuda_tests:
name: PEFT CUDA Tests
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
python -m pip install git+https://github.com/huggingface/peft.git
- name: Environment
run: |
python utils/print_env.py
- name: Run slow PEFT CUDA tests
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_peft_cuda \
tests/lora/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_peft_cuda_stats.txt
cat reports/tests_peft_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_peft_test_reports
path: reports
flax_tpu_tests:
name: Flax TPU Tests
runs-on: docker-tpu
container:
image: diffusers/diffusers-flax-tpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --privileged
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Run slow Flax TPU tests
if: ${{ matrix.config.framework == 'flax' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
--make-reports=tests_flax_tpu \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_flax_tpu_stats.txt
cat reports/tests_flax_tpu_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: flax_tpu_test_reports
path: reports
onnx_cuda_tests:
name: ONNX CUDA Tests
runs-on: docker-gpu
container:
image: diffusers/diffusers-onnxruntime-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Run slow ONNXRuntime CUDA tests
if: ${{ matrix.config.framework == 'onnxruntime' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_onnx_cuda \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_onnx_cuda_stats.txt
cat reports/tests_onnx_cuda_failures_short.txt
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: onnx_cuda_test_reports
path: reports
run_torch_compile_tests:
name: PyTorch Compile CUDA tests
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-compile-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m pip install -e .[quality,test,training]
- name: Environment
run: |
python utils/print_env.py
- name: Run example tests on GPU
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "compile" --make-reports=tests_torch_compile_cuda tests/
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_torch_compile_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_compile_test_reports
path: reports
run_xformers_tests:
name: PyTorch xformers CUDA tests
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-xformers-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m pip install -e .[quality,test,training]
- name: Environment
run: |
python utils/print_env.py
- name: Run example tests on GPU
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "xformers" --make-reports=tests_torch_xformers_cuda tests/
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_torch_xformers_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_xformers_test_reports
name: ${{ matrix.config.report }}_test_reports
path: reports
run_examples_tests:
@@ -415,13 +147,11 @@ jobs:
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/examples_torch_cuda_stats.txt
cat reports/examples_torch_cuda_failures_short.txt
run: cat reports/examples_torch_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: examples_test_reports
path: reports
path: reports

View File

@@ -40,7 +40,7 @@ jobs:
${CONDA_RUN} python -m pip install --upgrade pip
${CONDA_RUN} python -m pip install -e .[quality,test]
${CONDA_RUN} python -m pip install torch torchvision torchaudio
${CONDA_RUN} python -m pip install git+https://github.com/huggingface/accelerate.git
${CONDA_RUN} python -m pip install accelerate --upgrade
${CONDA_RUN} python -m pip install transformers --upgrade
- name: Environment

View File

@@ -17,7 +17,7 @@ jobs:
- name: Setup Python
uses: actions/setup-python@v1
with:
python-version: 3.8
python-version: 3.7
- name: Install requirements
run: |

View File

@@ -40,7 +40,7 @@ In the following, we give an overview of different ways to contribute, ranked by
As said before, **all contributions are valuable to the community**.
In the following, we will explain each contribution a bit more in detail.
For all contributions 4.-9. you will need to open a PR. It is explained in detail how to do so in [Opening a pull request](#how-to-open-a-pr)
For all contributions 4.-9. you will need to open a PR. It is explained in detail how to do so in [Opening a pull requst](#how-to-open-a-pr)
### 1. Asking and answering questions on the Diffusers discussion forum or on the Diffusers Discord
@@ -63,7 +63,7 @@ In the same spirit, you are of immense help to the community by answering such q
**Please** keep in mind that the more effort you put into asking or answering a question, the higher
the quality of the publicly documented knowledge. In the same way, well-posed and well-answered questions create a high-quality knowledge database accessible to everybody, while badly posed questions or answers reduce the overall quality of the public knowledge database.
In short, a high quality question or answer is *precise*, *concise*, *relevant*, *easy-to-understand*, *accessible*, and *well-formated/well-posed*. For more information, please have a look through the [How to write a good issue](#how-to-write-a-good-issue) section.
In short, a high quality question or answer is *precise*, *concise*, *relevant*, *easy-to-understand*, *accesible*, and *well-formated/well-posed*. For more information, please have a look through the [How to write a good issue](#how-to-write-a-good-issue) section.
**NOTE about channels**:
[*The forum*](https://discuss.huggingface.co/c/discussion-related-to-httpsgithubcomhuggingfacediffusers/63) is much better indexed by search engines, such as Google. Posts are ranked by popularity rather than chronologically. Hence, it's easier to look up questions and answers that we posted some time ago.
@@ -168,7 +168,7 @@ more precise, provide the link to a duplicated issue or redirect them to [the fo
If you have verified that the issued bug report is correct and requires a correction in the source code,
please have a look at the next sections.
For all of the following contributions, you will need to open a PR. It is explained in detail how to do so in the [Opening a pull request](#how-to-open-a-pr) section.
For all of the following contributions, you will need to open a PR. It is explained in detail how to do so in the [Opening a pull requst](#how-to-open-a-pr) section.
### 4. Fixing a "Good first issue"

View File

@@ -78,7 +78,7 @@ test:
# Run tests for examples
test-examples:
python -m pytest -n auto --dist=loadfile -s -v ./examples/
python -m pytest -n auto --dist=loadfile -s -v ./examples/pytorch/
# Release stuff

View File

@@ -70,7 +70,7 @@ The following design principles are followed:
- Pipelines should be used **only** for inference.
- Pipelines should be very readable, self-explanatory, and easy to tweak.
- Pipelines should be designed to build on top of each other and be easy to integrate into higher-level APIs.
- Pipelines are **not** intended to be feature-complete user interfaces. For future complete user interfaces one should rather have a look at [InvokeAI](https://github.com/invoke-ai/InvokeAI), [Diffuzers](https://github.com/abhishekkrthakur/diffuzers), and [lama-cleaner](https://github.com/Sanster/lama-cleaner).
- Pipelines are **not** intended to be feature-complete user interfaces. For future complete user interfaces one should rather have a look at [InvokeAI](https://github.com/invoke-ai/InvokeAI), [Diffuzers](https://github.com/abhishekkrthakur/diffuzers), and [lama-cleaner](https://github.com/Sanster/lama-cleaner)
- Every pipeline should have one and only one way to run it via a `__call__` method. The naming of the `__call__` arguments should be shared across all pipelines.
- Pipelines should be named after the task they are intended to solve.
- In almost all cases, novel diffusion pipelines shall be implemented in a new pipeline folder/file.
@@ -90,7 +90,7 @@ The following design principles are followed:
- To integrate new model checkpoints whose general architecture can be classified as an architecture that already exists in Diffusers, the existing model architecture shall be adapted to make it work with the new checkpoint. One should only create a new file if the model architecture is fundamentally different.
- Models should be designed to be easily extendable to future changes. This can be achieved by limiting public function arguments, configuration arguments, and "foreseeing" future changes, *e.g.* it is usually better to add `string` "...type" arguments that can easily be extended to new future types instead of boolean `is_..._type` arguments. Only the minimum amount of changes shall be made to existing architectures to make a new model checkpoint work.
- The model design is a difficult trade-off between keeping code readable and concise and supporting many model checkpoints. For most parts of the modeling code, classes shall be adapted for new model checkpoints, while there are some exceptions where it is preferred to add new classes to make sure the code is kept concise and
readable longterm, such as [UNet blocks](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_blocks.py) and [Attention processors](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention_processor.py).
readable longterm, such as [UNet blocks](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_blocks.py) and [Attention processors](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/cross_attention.py).
### Schedulers
@@ -104,7 +104,7 @@ The following design principles are followed:
- Schedulers all inherit from `SchedulerMixin` and `ConfigMixin`.
- Schedulers can be easily swapped out with the [`ConfigMixin.from_config`](https://huggingface.co/docs/diffusers/main/en/api/configuration#diffusers.ConfigMixin.from_config) method as explained in detail [here](./using-diffusers/schedulers.md).
- Every scheduler has to have a `set_num_inference_steps`, and a `step` function. `set_num_inference_steps(...)` has to be called before every denoising process, *i.e.* before `step(...)` is called.
- Every scheduler exposes the timesteps to be "looped over" via a `timesteps` attribute, which is an array of timesteps the model will be called upon.
- Every scheduler exposes the timesteps to be "looped over" via a `timesteps` attribute, which is an array of timesteps the model will be called upon
- The `step(...)` function takes a predicted model output and the "current" sample (x_t) and returns the "previous", slightly more denoised sample (x_t-1).
- Given the complexity of diffusion schedulers, the `step` function does not expose all the complexity and can be a bit of a "black box".
- In almost all cases, novel schedulers shall be implemented in a new scheduling file.

View File

@@ -1,6 +1,6 @@
<p align="center">
<br>
<img src="https://raw.githubusercontent.com/huggingface/diffusers/main/docs/source/en/imgs/diffusers_library.jpg" width="400"/>
<img src="https://github.com/huggingface/diffusers/blob/main/docs/source/en/imgs/diffusers_library.jpg" width="400"/>
<br>
<p>
<p align="center">
@@ -10,9 +10,6 @@
<a href="https://github.com/huggingface/diffusers/releases">
<img alt="GitHub release" src="https://img.shields.io/github/release/huggingface/diffusers.svg">
</a>
<a href="https://pepy.tech/project/diffusers">
<img alt="GitHub release" src="https://static.pepy.tech/badge/diffusers/month">
</a>
<a href="CODE_OF_CONDUCT.md">
<img alt="Contributor Covenant" src="https://img.shields.io/badge/Contributor%20Covenant-2.0-4baaaa.svg">
</a>

View File

@@ -1,46 +0,0 @@
FROM nvidia/cuda:12.1.0-runtime-ubuntu20.04
LABEL maintainer="Hugging Face"
LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.9 \
python3.9-dev \
python3-pip \
python3.9-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.9 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3.9 -m pip install --no-cache-dir --upgrade pip && \
python3.9 -m pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark && \
python3.9 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy \
scipy \
tensorboard \
transformers \
omegaconf
CMD ["/bin/bash"]

View File

@@ -1,4 +1,4 @@
FROM nvidia/cuda:12.1.0-runtime-ubuntu20.04
FROM nvidia/cuda:11.7.1-cudnn8-runtime-ubuntu20.04
LABEL maintainer="Hugging Face"
LABEL repository="diffusers"
@@ -6,16 +6,16 @@ ENV DEBIAN_FRONTEND=noninteractive
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.8 \
python3-pip \
python3.8-venv && \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.8 \
python3-pip \
python3.8-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
@@ -25,22 +25,23 @@ ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
python3 -m pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark && \
torch \
torchvision \
torchaudio \
invisible_watermark && \
python3 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy \
scipy \
tensorboard \
transformers \
omegaconf \
pytorch-lightning
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy \
scipy \
tensorboard \
transformers \
omegaconf \
pytorch-lightning \
xformers
CMD ["/bin/bash"]

View File

@@ -1,46 +0,0 @@
FROM nvidia/cuda:12.1.0-runtime-ubuntu20.04
LABEL maintainer="Hugging Face"
LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.8 \
python3-pip \
python3.8-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
python3 -m pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark && \
python3 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy \
scipy \
tensorboard \
transformers \
omegaconf \
xformers
CMD ["/bin/bash"]

View File

@@ -16,7 +16,7 @@ limitations under the License.
# Generating the documentation
To generate the documentation, you first have to build it. Several packages are necessary to build the doc,
To generate the documentation, you first have to build it. Several packages are necessary to build the doc,
you can install them with the following command, at the root of the code repository:
```bash
@@ -71,7 +71,7 @@ The `preview` command only works with existing doc files. When you add a complet
Accepted files are Markdown (.md).
Create a file with its extension and put it in the source directory. You can then link it to the toc-tree by putting
the filename without the extension in the [`_toctree.yml`](https://github.com/huggingface/diffusers/blob/main/docs/source/en/_toctree.yml) file.
the filename without the extension in the [`_toctree.yml`](https://github.com/huggingface/diffusers/blob/main/docs/source/_toctree.yml) file.
## Renaming section headers and moving sections
@@ -81,14 +81,14 @@ Therefore, we simply keep a little map of moved sections at the end of the docum
So if you renamed a section from: "Section A" to "Section B", then you can add at the end of the file:
```md
```
Sections that were moved:
[ <a href="#section-b">Section A</a><a id="section-a"></a> ]
```
and of course, if you moved it to another file, then:
```md
```
Sections that were moved:
[ <a href="../new-file#section-b">Section A</a><a id="section-a"></a> ]
@@ -109,8 +109,8 @@ although we can write them directly in Markdown.
Adding a new tutorial or section is done in two steps:
- Add a new Markdown (.md) file under `docs/source/<languageCode>`.
- Link that file in `docs/source/<languageCode>/_toctree.yml` on the correct toc-tree.
- Add a new file under `docs/source`. This file can either be ReStructuredText (.rst) or Markdown (.md).
- Link that file in `docs/source/_toctree.yml` on the correct toc-tree.
Make sure to put your new file under the proper section. It's unlikely to go in the first section (*Get Started*), so
depending on the intended targets (beginners, more advanced users, or researchers) it should go in sections two, three, or four.
@@ -119,7 +119,7 @@ depending on the intended targets (beginners, more advanced users, or researcher
When adding a new pipeline:
- Create a file `xxx.md` under `docs/source/<languageCode>/api/pipelines` (don't hesitate to copy an existing file as template).
- create a file `xxx.md` under `docs/source/api/pipelines` (don't hesitate to copy an existing file as template).
- Link that file in (*Diffusers Summary*) section in `docs/source/api/pipelines/overview.md`, along with the link to the paper, and a colab notebook (if available).
- Write a short overview of the diffusion model:
- Overview with paper & authors
@@ -129,6 +129,8 @@ When adding a new pipeline:
- Add all the pipeline classes that should be linked in the diffusion model. These classes should be added using our Markdown syntax. By default as follows:
```
## XXXPipeline
[[autodoc]] XXXPipeline
- all
- __call__
@@ -142,11 +144,11 @@ This will include every public method of the pipeline that is documented, as wel
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_xformers_memory_efficient_attention
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
```
You can follow the same process to create a new scheduler under the `docs/source/<languageCode>/api/schedulers` folder.
You can follow the same process to create a new scheduler under the `docs/source/api/schedulers` folder
### Writing source documentation
@@ -154,7 +156,7 @@ Values that should be put in `code` should either be surrounded by backticks: \`
and objects like True, None, or any strings should usually be put in `code`.
When mentioning a class, function, or method, it is recommended to use our syntax for internal links so that our tool
adds a link to its documentation with this syntax: \[\`XXXClass\`\] or \[\`function\`\]. This requires the class or
adds a link to its documentation with this syntax: \[\`XXXClass\`\] or \[\`function\`\]. This requires the class or
function to be in the main package.
If you want to create a link to some internal class or function, you need to
@@ -162,7 +164,7 @@ provide its path. For instance: \[\`pipelines.ImagePipelineOutput\`\]. This will
`pipelines.ImagePipelineOutput` in the description. To get rid of the path and only keep the name of the object you are
linking to in the description, add a ~: \[\`~pipelines.ImagePipelineOutput\`\] will generate a link with `ImagePipelineOutput` in the description.
The same works for methods so you can either use \[\`XXXClass.method\`\] or \[\`~XXXClass.method\`\].
The same works for methods so you can either use \[\`XXXClass.method\`\] or \[~\`XXXClass.method\`\].
#### Defining arguments in a method
@@ -194,8 +196,8 @@ Here's an example showcasing everything so far:
For optional arguments or arguments with defaults we follow the following syntax: imagine we have a function with the
following signature:
```py
def my_function(x: str=None, a: float=3.14):
```
def my_function(x: str = None, a: float = 1):
```
then its documentation should look like this:
@@ -204,7 +206,7 @@ then its documentation should look like this:
Args:
x (`str`, *optional*):
This argument controls ...
a (`float`, *optional*, defaults to `3.14`):
a (`float`, *optional*, defaults to 1):
This argument is used to ...
```
@@ -266,3 +268,4 @@ We have an automatic script running with the `make style` command that will make
This script may have some weird failures if you made a syntax mistake or if you uncover a bug. Therefore, it's
recommended to commit your changes before running `make style`, so you can revert the changes done by that script
easily.

View File

@@ -38,7 +38,7 @@ Here, `LANG-ID` should be one of the ISO 639-1 or ISO 639-2 language codes -- se
The fun part comes - translating the text!
The first thing we recommend is translating the part of the `_toctree.yml` file that corresponds to your doc chapter. This file is used to render the table of contents on the website.
The first thing we recommend is translating the part of the `_toctree.yml` file that corresponds to your doc chapter. This file is used to render the table of contents on the website.
> 🙋 If the `_toctree.yml` file doesn't yet exist for your language, you can create one by copy-pasting from the English version and deleting the sections unrelated to your chapter. Just make sure it exists in the `docs/source/LANG-ID/` directory!

View File

@@ -12,13 +12,9 @@
- local: tutorials/tutorial_overview
title: Overview
- local: using-diffusers/write_own_pipeline
title: Understanding pipelines, models and schedulers
- local: tutorials/autopipeline
title: AutoPipeline
title: Understanding models and schedulers
- local: tutorials/basic_training
title: Train a diffusion model
- local: tutorials/using_peft_for_inference
title: Inference with PEFT
title: Tutorials
- sections:
- sections:
@@ -29,15 +25,11 @@
- local: using-diffusers/schedulers
title: Load and compare different schedulers
- local: using-diffusers/custom_pipeline_overview
title: Load community pipelines and components
title: Load community pipelines
- local: using-diffusers/using_safetensors
title: Load safetensors
- local: using-diffusers/other-formats
title: Load different Stable Diffusion formats
- local: using-diffusers/loading_adapters
title: Load adapters
- local: using-diffusers/push_to_hub
title: Push files to the Hub
title: Loading & Hub
- sections:
- local: using-diffusers/pipeline_overview
@@ -45,15 +37,13 @@
- local: using-diffusers/unconditional_image_generation
title: Unconditional image generation
- local: using-diffusers/conditional_image_generation
title: Text-to-image
title: Text-to-image generation
- local: using-diffusers/img2img
title: Image-to-image
title: Text-guided image-to-image
- local: using-diffusers/inpaint
title: Inpainting
title: Text-guided image-inpainting
- local: using-diffusers/depth2img
title: Depth-to-image
title: Tasks
- sections:
title: Text-guided depth-to-image
- local: using-diffusers/textual_inversion_inference
title: Textual inversion
- local: training/distributed_inference
@@ -62,37 +52,17 @@
title: Improve image quality with deterministic generation
- local: using-diffusers/control_brightness
title: Control image brightness
- local: using-diffusers/weighted_prompts
title: Prompt weighting
- local: using-diffusers/freeu
title: Improve generation quality with FreeU
title: Techniques
- sections:
- local: using-diffusers/pipeline_overview
title: Overview
- local: using-diffusers/sdxl
title: Stable Diffusion XL
- local: using-diffusers/lcm
title: Latent Consistency Models
- local: using-diffusers/kandinsky
title: Kandinsky
- local: using-diffusers/controlnet
title: ControlNet
- local: using-diffusers/callback
title: Callback
- local: using-diffusers/shap-e
title: Shap-E
- local: using-diffusers/diffedit
title: DiffEdit
- local: using-diffusers/distilled_sd
title: Distilled Stable Diffusion inference
- local: using-diffusers/reproducibility
title: Create reproducible pipelines
- local: using-diffusers/custom_pipeline_examples
title: Community pipelines
- local: using-diffusers/contribute_pipeline
title: Contribute a community pipeline
title: Specific pipeline examples
title: How to contribute a community pipeline
- local: using-diffusers/stable_diffusion_jax_how_to
title: Stable Diffusion in JAX/Flax
- local: using-diffusers/weighted_prompts
title: Weighting Prompts
title: Pipelines for Inference
- sections:
- local: training/overview
title: Overview
@@ -116,10 +86,6 @@
title: InstructPix2Pix Training
- local: training/custom_diffusion
title: Custom Diffusion
- local: training/t2i_adapters
title: T2I-Adapters
- local: training/ddpo
title: Reinforcement learning training with DDPO
title: Training
- sections:
- local: using-diffusers/other-modalities
@@ -129,35 +95,25 @@
- sections:
- local: optimization/opt_overview
title: Overview
- sections:
- local: optimization/fp16
title: Speed up inference
- local: optimization/memory
title: Reduce memory usage
- local: optimization/torch2.0
title: PyTorch 2.0
- local: optimization/xformers
title: xFormers
- local: optimization/tome
title: Token merging
title: General optimizations
- sections:
- local: using-diffusers/stable_diffusion_jax_how_to
title: JAX/Flax
- local: optimization/onnx
title: ONNX
- local: optimization/open_vino
title: OpenVINO
- local: optimization/coreml
title: Core ML
title: Optimized model types
- sections:
- local: optimization/mps
title: Metal Performance Shaders (MPS)
- local: optimization/habana
title: Habana Gaudi
title: Optimized hardware
title: Optimization
- local: optimization/fp16
title: Memory and Speed
- local: optimization/torch2.0
title: Torch2.0 support
- local: optimization/xformers
title: xFormers
- local: optimization/onnx
title: ONNX
- local: optimization/open_vino
title: OpenVINO
- local: optimization/coreml
title: Core ML
- local: optimization/mps
title: MPS
- local: optimization/habana
title: Habana Gaudi
- local: optimization/tome
title: Token Merging
title: Optimization/Special Hardware
- sections:
- local: conceptual/philosophy
title: Philosophy
@@ -172,14 +128,22 @@
title: Conceptual Guides
- sections:
- sections:
- local: api/configuration
title: Configuration
- local: api/loaders
title: Loaders
- local: api/attnprocessor
title: Attention Processor
- local: api/diffusion_pipeline
title: Diffusion Pipeline
- local: api/logging
title: Logging
- local: api/configuration
title: Configuration
- local: api/outputs
title: Outputs
- local: api/loaders
title: Loaders
- local: api/utilities
title: Utilities
- local: api/image_processor
title: VAE Image Processor
title: Main Classes
- sections:
- local: api/models/overview
@@ -192,18 +156,12 @@
title: UNet2DConditionModel
- local: api/models/unet3d-cond
title: UNet3DConditionModel
- local: api/models/unet-motion
title: UNetMotionModel
- local: api/models/vq
title: VQModel
- local: api/models/autoencoderkl
title: AutoencoderKL
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_tiny
title: Tiny AutoEncoder
- local: api/models/consistency_decoder_vae
title: ConsistencyDecoderVAE
- local: api/models/transformer2d
title: Transformer2D
- local: api/models/transformer_temporal
@@ -218,26 +176,18 @@
title: Overview
- local: api/pipelines/alt_diffusion
title: AltDiffusion
- local: api/pipelines/animatediff
title: AnimateDiff
- local: api/pipelines/attend_and_excite
title: Attend-and-Excite
- local: api/pipelines/audio_diffusion
title: Audio Diffusion
- local: api/pipelines/audioldm
title: AudioLDM
- local: api/pipelines/audioldm2
title: AudioLDM 2
- local: api/pipelines/auto_pipeline
title: AutoPipeline
- local: api/pipelines/blip_diffusion
title: BLIP Diffusion
- local: api/pipelines/consistency_models
title: Consistency Models
- local: api/pipelines/controlnet
title: ControlNet
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/cycle_diffusion
title: Cycle Diffusion
- local: api/pipelines/dance_diffusion
@@ -255,25 +205,19 @@
- local: api/pipelines/pix2pix
title: InstructPix2Pix
- local: api/pipelines/kandinsky
title: Kandinsky 2.1
title: Kandinsky
- local: api/pipelines/kandinsky_v22
title: Kandinsky 2.2
- local: api/pipelines/latent_consistency_models
title: Latent Consistency Models
- local: api/pipelines/latent_diffusion
title: Latent Diffusion
- local: api/pipelines/panorama
title: MultiDiffusion
- local: api/pipelines/musicldm
title: MusicLDM
- local: api/pipelines/paint_by_example
title: Paint By Example
title: PaintByExample
- local: api/pipelines/paradigms
title: Parallel Sampling of Diffusion Models
- local: api/pipelines/pix2pix_zero
title: Pix2Pix Zero
- local: api/pipelines/pixart
title: PixArt
- local: api/pipelines/pndm
title: PNDM
- local: api/pipelines/repaint
@@ -314,9 +258,7 @@
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
title: LDM3D Text-to-(RGB, Depth)
- local: api/pipelines/stable_diffusion/adapter
title: Stable Diffusion T2I-Adapter
- local: api/pipelines/stable_diffusion/gligen
title: GLIGEN (Grounded Language-to-Image Generation)
title: Stable Diffusion T2I-adapter
title: Stable Diffusion
- local: api/pipelines/stable_unclip
title: Stable unCLIP
@@ -329,7 +271,7 @@
- local: api/pipelines/text_to_video_zero
title: Text2Video-Zero
- local: api/pipelines/unclip
title: unCLIP
title: UnCLIP
- local: api/pipelines/latent_diffusion_uncond
title: Unconditional Latent Diffusion
- local: api/pipelines/unidiffuser
@@ -340,75 +282,55 @@
title: Versatile Diffusion
- local: api/pipelines/vq_diffusion
title: VQ Diffusion
- local: api/pipelines/wuerstchen
title: Wuerstchen
title: Pipelines
- sections:
- local: api/schedulers/overview
title: Overview
- local: api/schedulers/cm_stochastic_iterative
title: CMStochasticIterativeScheduler
- local: api/schedulers/consistency_decoder
title: ConsistencyDecoderScheduler
- local: api/schedulers/ddim_inverse
title: DDIMInverseScheduler
title: Consistency Model Multistep Scheduler
- local: api/schedulers/ddim
title: DDIMScheduler
title: DDIM
- local: api/schedulers/ddim_inverse
title: DDIMInverse
- local: api/schedulers/ddpm
title: DDPMScheduler
title: DDPM
- local: api/schedulers/deis
title: DEISMultistepScheduler
- local: api/schedulers/multistep_dpm_solver_inverse
title: DPMSolverMultistepInverse
- local: api/schedulers/multistep_dpm_solver
title: DPMSolverMultistepScheduler
title: DEIS
- local: api/schedulers/dpm_discrete
title: DPM Discrete Scheduler
- local: api/schedulers/dpm_discrete_ancestral
title: DPM Discrete Scheduler with ancestral sampling
- local: api/schedulers/dpm_sde
title: DPMSolverSDEScheduler
- local: api/schedulers/singlestep_dpm_solver
title: DPMSolverSinglestepScheduler
- local: api/schedulers/euler_ancestral
title: EulerAncestralDiscreteScheduler
title: Euler Ancestral Scheduler
- local: api/schedulers/euler
title: EulerDiscreteScheduler
title: Euler scheduler
- local: api/schedulers/heun
title: HeunDiscreteScheduler
title: Heun Scheduler
- local: api/schedulers/multistep_dpm_solver_inverse
title: Inverse Multistep DPM-Solver
- local: api/schedulers/ipndm
title: IPNDMScheduler
- local: api/schedulers/stochastic_karras_ve
title: KarrasVeScheduler
- local: api/schedulers/dpm_discrete_ancestral
title: KDPM2AncestralDiscreteScheduler
- local: api/schedulers/dpm_discrete
title: KDPM2DiscreteScheduler
- local: api/schedulers/lcm
title: LCMScheduler
title: IPNDM
- local: api/schedulers/lms_discrete
title: LMSDiscreteScheduler
title: Linear Multistep
- local: api/schedulers/multistep_dpm_solver
title: Multistep DPM-Solver
- local: api/schedulers/pndm
title: PNDMScheduler
title: PNDM
- local: api/schedulers/repaint
title: RePaintScheduler
- local: api/schedulers/score_sde_ve
title: ScoreSdeVeScheduler
- local: api/schedulers/score_sde_vp
title: ScoreSdeVpScheduler
title: RePaint Scheduler
- local: api/schedulers/singlestep_dpm_solver
title: Singlestep DPM-Solver
- local: api/schedulers/stochastic_karras_ve
title: Stochastic Kerras VE
- local: api/schedulers/unipc
title: UniPCMultistepScheduler
- local: api/schedulers/score_sde_ve
title: VE-SDE
- local: api/schedulers/score_sde_vp
title: VP-SDE
- local: api/schedulers/vq_diffusion
title: VQDiffusionScheduler
title: Schedulers
- sections:
- local: api/internal_classes_overview
title: Overview
- local: api/attnprocessor
title: Attention Processor
- local: api/activations
title: Custom activation functions
- local: api/normalization
title: Custom normalization layers
- local: api/utilities
title: Utilities
- local: api/image_processor
title: VAE Image Processor
title: Internal classes
title: API

View File

@@ -1,15 +0,0 @@
# Activation functions
Customized activation functions for supporting various models in 🤗 Diffusers.
## GELU
[[autodoc]] models.activations.GELU
## GEGLU
[[autodoc]] models.activations.GEGLU
## ApproximateGELU
[[autodoc]] models.activations.ApproximateGELU

View File

@@ -17,9 +17,6 @@ An attention processor is a class for applying different types of attention mech
## CustomDiffusionAttnProcessor
[[autodoc]] models.attention_processor.CustomDiffusionAttnProcessor
## CustomDiffusionAttnProcessor2_0
[[autodoc]] models.attention_processor.CustomDiffusionAttnProcessor2_0
## AttnAddedKVProcessor
[[autodoc]] models.attention_processor.AttnAddedKVProcessor
@@ -42,4 +39,4 @@ An attention processor is a class for applying different types of attention mech
[[autodoc]] models.attention_processor.SlicedAttnProcessor
## SlicedAttnAddedKVProcessor
[[autodoc]] models.attention_processor.SlicedAttnAddedKVProcessor
[[autodoc]] models.attention_processor.SlicedAttnAddedKVProcessor

View File

@@ -0,0 +1,36 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Pipelines
The [`DiffusionPipeline`] is the quickest way to load any pretrained diffusion pipeline from the [Hub](https://huggingface.co/models?library=diffusers) for inference.
<Tip>
You shouldn't use the [`DiffusionPipeline`] class for training or finetuning a diffusion model. Individual
components (for example, [`UNet2DModel`] and [`UNet2DConditionModel`]) of diffusion pipelines are usually trained individually, so we suggest directly working with them instead.
</Tip>
The pipeline type (for example [`StableDiffusionPipeline`]) of any diffusion pipeline loaded with [`~DiffusionPipeline.from_pretrained`] is automatically
detected and pipeline components are loaded and passed to the `__init__` function of the pipeline.
Any pipeline object can be saved locally with [`~DiffusionPipeline.save_pretrained`].
## DiffusionPipeline
[[autodoc]] DiffusionPipeline
- all
- __call__
- device
- to
- components

View File

@@ -1,3 +0,0 @@
# Overview
The APIs in this section are more experimental and prone to breaking changes. Most of them are used internally for development, but they may also be useful to you if you're interested in building a diffusion model with some custom parts or if you're interested in some of our helper utilities for working with 🤗 Diffusers.

View File

@@ -28,10 +28,6 @@ Adapters (textual inversion, LoRA, hypernetworks) allow you to modify a diffusio
[[autodoc]] loaders.TextualInversionLoaderMixin
## StableDiffusionXLLoraLoaderMixin
[[autodoc]] loaders.StableDiffusionXLLoraLoaderMixin
## LoraLoaderMixin
[[autodoc]] loaders.LoraLoaderMixin

View File

@@ -67,30 +67,30 @@ By default, `tqdm` progress bars are displayed during model download. [`logging.
## Base setters
[[autodoc]] utils.logging.set_verbosity_error
[[autodoc]] logging.set_verbosity_error
[[autodoc]] utils.logging.set_verbosity_warning
[[autodoc]] logging.set_verbosity_warning
[[autodoc]] utils.logging.set_verbosity_info
[[autodoc]] logging.set_verbosity_info
[[autodoc]] utils.logging.set_verbosity_debug
[[autodoc]] logging.set_verbosity_debug
## Other functions
[[autodoc]] utils.logging.get_verbosity
[[autodoc]] logging.get_verbosity
[[autodoc]] utils.logging.set_verbosity
[[autodoc]] logging.set_verbosity
[[autodoc]] utils.logging.get_logger
[[autodoc]] logging.get_logger
[[autodoc]] utils.logging.enable_default_handler
[[autodoc]] logging.enable_default_handler
[[autodoc]] utils.logging.disable_default_handler
[[autodoc]] logging.disable_default_handler
[[autodoc]] utils.logging.enable_explicit_format
[[autodoc]] logging.enable_explicit_format
[[autodoc]] utils.logging.reset_format
[[autodoc]] logging.reset_format
[[autodoc]] utils.logging.enable_progress_bar
[[autodoc]] logging.enable_progress_bar
[[autodoc]] utils.logging.disable_progress_bar
[[autodoc]] logging.disable_progress_bar

View File

@@ -1,45 +0,0 @@
# Tiny AutoEncoder
Tiny AutoEncoder for Stable Diffusion (TAESD) was introduced in [madebyollin/taesd](https://github.com/madebyollin/taesd) by Ollin Boer Bohan. It is a tiny distilled version of Stable Diffusion's VAE that can quickly decode the latents in a [`StableDiffusionPipeline`] or [`StableDiffusionXLPipeline`] almost instantly.
To use with Stable Diffusion v-2.1:
```python
import torch
from diffusers import DiffusionPipeline, AutoencoderTiny
pipe = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-1-base", torch_dtype=torch.float16
)
pipe.vae = AutoencoderTiny.from_pretrained("madebyollin/taesd", torch_dtype=torch.float16)
pipe = pipe.to("cuda")
prompt = "slice of delicious New York-style berry cheesecake"
image = pipe(prompt, num_inference_steps=25).images[0]
image.save("cheesecake.png")
```
To use with Stable Diffusion XL 1.0
```python
import torch
from diffusers import DiffusionPipeline, AutoencoderTiny
pipe = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16
)
pipe.vae = AutoencoderTiny.from_pretrained("madebyollin/taesdxl", torch_dtype=torch.float16)
pipe = pipe.to("cuda")
prompt = "slice of delicious New York-style berry cheesecake"
image = pipe(prompt, num_inference_steps=25).images[0]
image.save("cheesecake_sdxl.png")
```
## AutoencoderTiny
[[autodoc]] AutoencoderTiny
## AutoencoderTinyOutput
[[autodoc]] models.autoencoder_tiny.AutoencoderTinyOutput

View File

@@ -1,18 +0,0 @@
# Consistency Decoder
Consistency decoder can be used to decode the latents from the denoising UNet in the [`StableDiffusionPipeline`]. This decoder was introduced in the [DALL-E 3 technical report](https://openai.com/dall-e-3).
The original codebase can be found at [openai/consistencydecoder](https://github.com/openai/consistencydecoder).
<Tip warning={true}>
Inference is only supported for 2 iterations as of now.
</Tip>
The pipeline could not have been contributed without the help of [madebyollin](https://github.com/madebyollin) and [mrsteyk](https://github.com/mrsteyk) from [this issue](https://github.com/openai/consistencydecoder/issues/1).
## ConsistencyDecoderVAE
[[autodoc]] ConsistencyDecoderVAE
- all
- decode

View File

@@ -12,13 +12,13 @@ By default the [`ControlNetModel`] should be loaded with [`~ModelMixin.from_pret
from the original format using [`FromOriginalControlnetMixin.from_single_file`] as follows:
```py
from diffusers import StableDiffusionControlNetPipeline, ControlNetModel
from diffusers import StableDiffusionControlnetPipeline, ControlNetModel
url = "https://huggingface.co/lllyasviel/ControlNet-v1-1/blob/main/control_v11p_sd15_canny.pth" # can also be a local path
controlnet = ControlNetModel.from_single_file(url)
url = "https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/v1-5-pruned.safetensors" # can also be a local path
pipe = StableDiffusionControlNetPipeline.from_single_file(url, controlnet=controlnet)
pipe = StableDiffusionControlnetPipeline.from_single_file(url, controlnet=controlnet)
```
## ControlNetModel

View File

@@ -9,8 +9,4 @@ All models are built from the base [`ModelMixin`] class which is a [`torch.nn.mo
## FlaxModelMixin
[[autodoc]] FlaxModelMixin
## PushToHubMixin
[[autodoc]] utils.PushToHubMixin
[[autodoc]] FlaxModelMixin

View File

@@ -1,13 +0,0 @@
# UNetMotionModel
The [UNet](https://huggingface.co/papers/1505.04597) model was originally introduced by Ronneberger et al for biomedical image segmentation, but it is also commonly used in 🤗 Diffusers because it outputs images that are the same size as the input. It is one of the most important components of a diffusion system because it facilitates the actual diffusion process. There are several variants of the UNet model in 🤗 Diffusers, depending on it's number of dimensions and whether it is a conditional model or not. This is a 2D UNet model.
The abstract from the paper is:
*There is large consent that successful training of deep networks requires many thousand annotated training samples. In this paper, we present a network and training strategy that relies on the strong use of data augmentation to use the available annotated samples more efficiently. The architecture consists of a contracting path to capture context and a symmetric expanding path that enables precise localization. We show that such a network can be trained end-to-end from very few images and outperforms the prior best method (a sliding-window convolutional network) on the ISBI challenge for segmentation of neuronal structures in electron microscopic stacks. Using the same network trained on transmitted light microscopy images (phase contrast and DIC) we won the ISBI cell tracking challenge 2015 in these categories by a large margin. Moreover, the network is fast. Segmentation of a 512x512 image takes less than a second on a recent GPU. The full implementation (based on Caffe) and the trained networks are available at http://lmb.informatik.uni-freiburg.de/people/ronneber/u-net.*
## UNetMotionModel
[[autodoc]] UNetMotionModel
## UNet3DConditionOutput
[[autodoc]] models.unet_3d_condition.UNet3DConditionOutput

View File

@@ -1,15 +0,0 @@
# Normalization layers
Customized normalization layers for supporting various models in 🤗 Diffusers.
## AdaLayerNorm
[[autodoc]] models.normalization.AdaLayerNorm
## AdaLayerNormZero
[[autodoc]] models.normalization.AdaLayerNormZero
## AdaGroupNorm
[[autodoc]] models.normalization.AdaGroupNorm

View File

@@ -24,7 +24,7 @@ The abstract from the paper is:
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -1,230 +0,0 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Text-to-Video Generation with AnimateDiff
## Overview
[AnimateDiff: Animate Your Personalized Text-to-Image Diffusion Models without Specific Tuning](https://arxiv.org/abs/2307.04725) by Yuwei Guo, Ceyuan Yang*, Anyi Rao, Yaohui Wang, Yu Qiao, Dahua Lin, Bo Dai
The abstract of the paper is the following:
With the advance of text-to-image models (e.g., Stable Diffusion) and corresponding personalization techniques such as DreamBooth and LoRA, everyone can manifest their imagination into high-quality images at an affordable cost. Subsequently, there is a great demand for image animation techniques to further combine generated static images with motion dynamics. In this report, we propose a practical framework to animate most of the existing personalized text-to-image models once and for all, saving efforts in model-specific tuning. At the core of the proposed framework is to insert a newly initialized motion modeling module into the frozen text-to-image model and train it on video clips to distill reasonable motion priors. Once trained, by simply injecting this motion modeling module, all personalized versions derived from the same base T2I readily become text-driven models that produce diverse and personalized animated images. We conduct our evaluation on several public representative personalized text-to-image models across anime pictures and realistic photographs, and demonstrate that our proposed framework helps these models generate temporally smooth animation clips while preserving the domain and diversity of their outputs. Code and pre-trained weights will be publicly available at this https URL .
## Available Pipelines
| Pipeline | Tasks | Demo
|---|---|:---:|
| [AnimateDiffPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff.py) | *Text-to-Video Generation with AnimateDiff* |
## Available checkpoints
Motion Adapter checkpoints can be found under [guoyww](https://huggingface.co/guoyww/). These checkpoints are meant to work with any model based on Stable Diffusion 1.4/1.5
## Usage example
AnimateDiff works with a MotionAdapter checkpoint and a Stable Diffusion model checkpoint. The MotionAdapter is a collection of Motion Modules that are responsible for adding coherent motion across image frames. These modules are applied after the Resnet and Attention blocks in Stable Diffusion UNet.
The following example demonstrates how to use a *MotionAdapter* checkpoint with Diffusers for inference based on StableDiffusion-1.4/1.5.
```python
import torch
from diffusers import MotionAdapter, AnimateDiffPipeline, DDIMScheduler
from diffusers.utils import export_to_gif
# Load the motion adapter
adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-v1-5-2")
# load SD 1.5 based finetuned model
model_id = "SG161222/Realistic_Vision_V5.1_noVAE"
pipe = AnimateDiffPipeline.from_pretrained(model_id, motion_adapter=adapter)
scheduler = DDIMScheduler.from_pretrained(
model_id, subfolder="scheduler", clip_sample=False, timestep_spacing="linspace", steps_offset=1
)
pipe.scheduler = scheduler
# enable memory savings
pipe.enable_vae_slicing()
pipe.enable_model_cpu_offload()
output = pipe(
prompt=(
"masterpiece, bestquality, highlydetailed, ultradetailed, sunset, "
"orange sky, warm lighting, fishing boats, ocean waves seagulls, "
"rippling water, wharf, silhouette, serene atmosphere, dusk, evening glow, "
"golden hour, coastal landscape, seaside scenery"
),
negative_prompt="bad quality, worse quality",
num_frames=16,
guidance_scale=7.5,
num_inference_steps=25,
generator=torch.Generator("cpu").manual_seed(42),
)
frames = output.frames[0]
export_to_gif(frames, "animation.gif")
```
Here are some sample outputs:
<table>
<tr>
<td><center>
masterpiece, bestquality, sunset.
<br>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-realistic-doc.gif"
alt="masterpiece, bestquality, sunset"
style="width: 300px;" />
</center></td>
</tr>
</table>
<Tip>
AnimateDiff tends to work better with finetuned Stable Diffusion models. If you plan on using a scheduler that can clip samples, make sure to disable it by setting `clip_sample=False` in the scheduler as this can also have an adverse effect on generated samples.
</Tip>
## Using Motion LoRAs
Motion LoRAs are a collection of LoRAs that work with the `guoyww/animatediff-motion-adapter-v1-5-2` checkpoint. These LoRAs are responsible for adding specific types of motion to the animations.
```python
import torch
from diffusers import MotionAdapter, AnimateDiffPipeline, DDIMScheduler
from diffusers.utils import export_to_gif
# Load the motion adapter
adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-v1-5-2")
# load SD 1.5 based finetuned model
model_id = "SG161222/Realistic_Vision_V5.1_noVAE"
pipe = AnimateDiffPipeline.from_pretrained(model_id, motion_adapter=adapter)
pipe.load_lora_weights("guoyww/animatediff-motion-lora-zoom-out", adapter_name="zoom-out")
scheduler = DDIMScheduler.from_pretrained(
model_id, subfolder="scheduler", clip_sample=False, timestep_spacing="linspace", steps_offset=1
)
pipe.scheduler = scheduler
# enable memory savings
pipe.enable_vae_slicing()
pipe.enable_model_cpu_offload()
output = pipe(
prompt=(
"masterpiece, bestquality, highlydetailed, ultradetailed, sunset, "
"orange sky, warm lighting, fishing boats, ocean waves seagulls, "
"rippling water, wharf, silhouette, serene atmosphere, dusk, evening glow, "
"golden hour, coastal landscape, seaside scenery"
),
negative_prompt="bad quality, worse quality",
num_frames=16,
guidance_scale=7.5,
num_inference_steps=25,
generator=torch.Generator("cpu").manual_seed(42),
)
frames = output.frames[0]
export_to_gif(frames, "animation.gif")
```
<table>
<tr>
<td><center>
masterpiece, bestquality, sunset.
<br>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-zoom-out-lora.gif"
alt="masterpiece, bestquality, sunset"
style="width: 300px;" />
</center></td>
</tr>
</table>
## Using Motion LoRAs with PEFT
You can also leverage the [PEFT](https://github.com/huggingface/peft) backend to combine Motion LoRA's and create more complex animations.
First install PEFT with
```shell
pip install peft
```
Then you can use the following code to combine Motion LoRAs.
```python
import torch
from diffusers import MotionAdapter, AnimateDiffPipeline, DDIMScheduler
from diffusers.utils import export_to_gif
# Load the motion adapter
adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-v1-5-2")
# load SD 1.5 based finetuned model
model_id = "SG161222/Realistic_Vision_V5.1_noVAE"
pipe = AnimateDiffPipeline.from_pretrained(model_id, motion_adapter=adapter)
pipe.load_lora_weights("diffusers/animatediff-motion-lora-zoom-out", adapter_name="zoom-out")
pipe.load_lora_weights("diffusers/animatediff-motion-lora-pan-left", adapter_name="pan-left")
pipe.set_adapters(["zoom-out", "pan-left"], adapter_weights=[1.0, 1.0])
scheduler = DDIMScheduler.from_pretrained(
model_id, subfolder="scheduler", clip_sample=False, timestep_spacing="linspace", steps_offset=1
)
pipe.scheduler = scheduler
# enable memory savings
pipe.enable_vae_slicing()
pipe.enable_model_cpu_offload()
output = pipe(
prompt=(
"masterpiece, bestquality, highlydetailed, ultradetailed, sunset, "
"orange sky, warm lighting, fishing boats, ocean waves seagulls, "
"rippling water, wharf, silhouette, serene atmosphere, dusk, evening glow, "
"golden hour, coastal landscape, seaside scenery"
),
negative_prompt="bad quality, worse quality",
num_frames=16,
guidance_scale=7.5,
num_inference_steps=25,
generator=torch.Generator("cpu").manual_seed(42),
)
frames = output.frames[0]
export_to_gif(frames, "animation.gif")
```
<table>
<tr>
<td><center>
masterpiece, bestquality, sunset.
<br>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-zoom-out-pan-left-lora.gif"
alt="masterpiece, bestquality, sunset"
style="width: 300px;" />
</center></td>
</tr>
</table>
## AnimateDiffPipeline
[[autodoc]] AnimateDiffPipeline
- all
- __call__
- enable_freeu
- disable_freeu
- enable_vae_slicing
- disable_vae_slicing
- enable_vae_tiling
- disable_vae_tiling
## AnimateDiffPipelineOutput
[[autodoc]] pipelines.animatediff.AnimateDiffPipelineOutput

View File

@@ -22,7 +22,7 @@ You can find additional information about Attend-and-Excite on the [project page
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -18,7 +18,7 @@ The original codebase, training scripts and example notebooks can be found at [t
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -37,7 +37,7 @@ During inference:
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
@@ -46,5 +46,6 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
- all
- __call__
## AudioPipelineOutput
[[autodoc]] pipelines.AudioPipelineOutput
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput

View File

@@ -1,91 +0,0 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# AudioLDM 2
AudioLDM 2 was proposed in [AudioLDM 2: Learning Holistic Audio Generation with Self-supervised Pretraining](https://arxiv.org/abs/2308.05734)
by Haohe Liu et al. AudioLDM 2 takes a text prompt as input and predicts the corresponding audio. It can generate
text-conditional sound effects, human speech and music.
Inspired by [Stable Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/overview), AudioLDM 2
is a text-to-audio _latent diffusion model (LDM)_ that learns continuous audio representations from text embeddings. Two
text encoder models are used to compute the text embeddings from a prompt input: the text-branch of [CLAP](https://huggingface.co/docs/transformers/main/en/model_doc/clap)
and the encoder of [Flan-T5](https://huggingface.co/docs/transformers/main/en/model_doc/flan-t5). These text embeddings
are then projected to a shared embedding space by an [AudioLDM2ProjectionModel](https://huggingface.co/docs/diffusers/main/api/pipelines/audioldm2#diffusers.AudioLDM2ProjectionModel).
A [GPT2](https://huggingface.co/docs/transformers/main/en/model_doc/gpt2) _language model (LM)_ is used to auto-regressively
predict eight new embedding vectors, conditional on the projected CLAP and Flan-T5 embeddings. The generated embedding
vectors and Flan-T5 text embeddings are used as cross-attention conditioning in the LDM. The [UNet](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2UNet2DConditionModel)
of AudioLDM 2 is unique in the sense that it takes **two** cross-attention embeddings, as opposed to one cross-attention
conditioning, as in most other LDMs.
The abstract of the paper is the following:
*Although audio generation shares commonalities across different types of audio, such as speech, music, and sound effects, designing models for each type requires careful consideration of specific objectives and biases that can significantly differ from those of other types. To bring us closer to a unified perspective of audio generation, this paper proposes a framework that utilizes the same learning method for speech, music, and sound effect generation. Our framework introduces a general representation of audio, called language of audio (LOA). Any audio can be translated into LOA based on AudioMAE, a self-supervised pre-trained representation learning model. In the generation process, we translate any modalities into LOA by using a GPT-2 model, and we perform self-supervised audio generation learning with a latent diffusion model conditioned on LOA. The proposed framework naturally brings advantages such as in-context learning abilities and reusable self-supervised pretrained AudioMAE and latent diffusion models. Experiments on the major benchmarks of text-to-audio, text-to-music, and text-to-speech demonstrate new state-of-the-art or competitive performance to previous approaches.*
This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi). The original codebase can be
found at [haoheliu/audioldm2](https://github.com/haoheliu/audioldm2).
## Tips
### Choosing a checkpoint
AudioLDM2 comes in three variants. Two of these checkpoints are applicable to the general task of text-to-audio
generation. The third checkpoint is trained exclusively on text-to-music generation.
All checkpoints share the same model size for the text encoders and VAE. They differ in the size and depth of the UNet.
See table below for details on the three checkpoints:
| Checkpoint | Task | UNet Model Size | Total Model Size | Training Data / h |
|-----------------------------------------------------------------|---------------|-----------------|------------------|-------------------|
| [audioldm2](https://huggingface.co/cvssp/audioldm2) | Text-to-audio | 350M | 1.1B | 1150k |
| [audioldm2-large](https://huggingface.co/cvssp/audioldm2-large) | Text-to-audio | 750M | 1.5B | 1150k |
| [audioldm2-music](https://huggingface.co/cvssp/audioldm2-music) | Text-to-music | 350M | 1.1B | 665k |
### Constructing a prompt
* Descriptive prompt inputs work best: use adjectives to describe the sound (e.g. "high quality" or "clear") and make the prompt context specific (e.g. "water stream in a forest" instead of "stream").
* It's best to use general terms like "cat" or "dog" instead of specific names or abstract objects the model may not be familiar with.
* Using a **negative prompt** can significantly improve the quality of the generated waveform, by guiding the generation away from terms that correspond to poor quality audio. Try using a negative prompt of "Low quality."
### Controlling inference
* The _quality_ of the predicted audio sample can be controlled by the `num_inference_steps` argument; higher steps give higher quality audio at the expense of slower inference.
* The _length_ of the predicted audio sample can be controlled by varying the `audio_length_in_s` argument.
### Evaluating generated waveforms:
* The quality of the generated waveforms can vary significantly based on the seed. Try generating with different seeds until you find a satisfactory generation
* Multiple waveforms can be generated in one go: set `num_waveforms_per_prompt` to a value greater than 1. Automatic scoring will be performed between the generated waveforms and prompt text, and the audios ranked from best to worst accordingly.
The following example demonstrates how to construct good music generation using the aforementioned tips: [example](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2Pipeline.__call__.example).
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## AudioLDM2Pipeline
[[autodoc]] AudioLDM2Pipeline
- all
- __call__
## AudioLDM2ProjectionModel
[[autodoc]] AudioLDM2ProjectionModel
- forward
## AudioLDM2UNet2DConditionModel
[[autodoc]] AudioLDM2UNet2DConditionModel
- forward
## AudioPipelineOutput
[[autodoc]] pipelines.AudioPipelineOutput

View File

@@ -12,41 +12,35 @@ specific language governing permissions and limitations under the License.
# AutoPipeline
`AutoPipeline` is designed to:
In many cases, one checkpoint can be used for multiple tasks. For example, you may be able to use the same checkpoint for Text-to-Image, Image-to-Image, and Inpainting. However, you'll need to know the pipeline class names linked to your checkpoint.
1. make it easy for you to load a checkpoint for a task without knowing the specific pipeline class to use
2. use multiple pipelines in your workflow
AutoPipeline is designed to make it easy for you to use multiple pipelines in your workflow. We currently provide 3 AutoPipeline classes to perform three different tasks, i.e. [`AutoPipelineForText2Image`], [`AutoPipelineForImage2Image`], and [`AutoPipelineForInpainting`]. You'll need to choose the AutoPipeline class based on the task you want to perform and use it to automatically retrieve the relevant pipeline given the name/path to the pre-trained weights.
Based on the task, the `AutoPipeline` class automatically retrieves the relevant pipeline given the name or path to the pretrained weights with the `from_pretrained()` method.
For example, to perform Image-to-Image with the SD1.5 checkpoint, you can do
To seamlessly switch between tasks with the same checkpoint without reallocating additional memory, use the `from_pipe()` method to transfer the components from the original pipeline to the new one.
```python
from diffusers import PipelineForImageToImage
```py
from diffusers import AutoPipelineForText2Image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
image = pipeline(prompt, num_inference_steps=25).images[0]
pipe_i2i = PipelineForImageoImage.from_pretrained("runwayml/stable-diffusion-v1-5")
```
<Tip>
It will also help you switch between tasks seamlessly using the same checkpoint without reallocating additional memory. For example, to re-use the Image-to-Image pipeline we just created for inpainting, you can do
Check out the [AutoPipeline](/tutorials/autopipeline) tutorial to learn how to use this API!
```python
from diffusers import PipelineForInpainting
</Tip>
pipe_inpaint = AutoPipelineForInpainting.from_pipe(pipe_i2i)
```
All the components will be transferred to the inpainting pipeline with zero cost.
`AutoPipeline` supports text-to-image, image-to-image, and inpainting for the following diffusion models:
- [Stable Diffusion](./stable_diffusion)
- [ControlNet](./controlnet)
- [Stable Diffusion XL (SDXL)](./stable_diffusion/stable_diffusion_xl)
- [DeepFloyd IF](./if)
Currently AutoPipeline support the Text-to-Image, Image-to-Image, and Inpainting tasks for below diffusion models:
- [stable Diffusion](./stable_diffusion)
- [Stable Diffusion Controlnet](./api/pipelines/controlnet)
- [Stable Diffusion XL](./stable_diffusion/stable_diffusion_xl)
- [IF](./if)
- [Kandinsky](./kandinsky)
- [Kandinsky 2.2](./kandinsky#kandinsky-22)
- [Kandinsky 2.2](./kandinsky)
## AutoPipelineForText2Image

View File

@@ -1,29 +0,0 @@
# Blip Diffusion
Blip Diffusion was proposed in [BLIP-Diffusion: Pre-trained Subject Representation for Controllable Text-to-Image Generation and Editing](https://arxiv.org/abs/2305.14720). It enables zero-shot subject-driven generation and control-guided zero-shot generation.
The abstract from the paper is:
*Subject-driven text-to-image generation models create novel renditions of an input subject based on text prompts. Existing models suffer from lengthy fine-tuning and difficulties preserving the subject fidelity. To overcome these limitations, we introduce BLIP-Diffusion, a new subject-driven image generation model that supports multimodal control which consumes inputs of subject images and text prompts. Unlike other subject-driven generation models, BLIP-Diffusion introduces a new multimodal encoder which is pre-trained to provide subject representation. We first pre-train the multimodal encoder following BLIP-2 to produce visual representation aligned with the text. Then we design a subject representation learning task which enables a diffusion model to leverage such visual representation and generates new subject renditions. Compared with previous methods such as DreamBooth, our model enables zero-shot subject-driven generation, and efficient fine-tuning for customized subject with up to 20x speedup. We also demonstrate that BLIP-Diffusion can be flexibly combined with existing techniques such as ControlNet and prompt-to-prompt to enable novel subject-driven generation and editing applications.*
The original codebase can be found at [salesforce/LAVIS](https://github.com/salesforce/LAVIS/tree/main/projects/blip-diffusion). You can find the official BLIP Diffusion checkpoints under the [hf.co/SalesForce](https://hf.co/SalesForce) organization.
`BlipDiffusionPipeline` and `BlipDiffusionControlNetPipeline` were contributed by [`ayushtues`](https://github.com/ayushtues/).
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## BlipDiffusionPipeline
[[autodoc]] BlipDiffusionPipeline
- all
- __call__
## BlipDiffusionControlNetPipeline
[[autodoc]] BlipDiffusionControlNetPipeline
- all
- __call__

View File

@@ -12,9 +12,9 @@ specific language governing permissions and limitations under the License.
# ControlNet
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang and Maneesh Agrawala.
[Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
Using a pretrained model, we can provide control images (for example, a depth map) to control Stable Diffusion text-to-image generation so that it follows the structure of the depth image and fills in the details.
The abstract from the paper is:
@@ -22,13 +22,290 @@ The abstract from the paper is:
This model was contributed by [takuma104](https://huggingface.co/takuma104). ❤️
The original codebase can be found at [lllyasviel/ControlNet](https://github.com/lllyasviel/ControlNet), and you can find official ControlNet checkpoints on [lllyasviel's](https://huggingface.co/lllyasviel) Hub profile.
The original codebase can be found at [lllyasviel/ControlNet](https://github.com/lllyasviel/ControlNet).
<Tip>
## Usage example
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
In the following we give a simple example of how to use a *ControlNet* checkpoint with Diffusers for inference.
The inference pipeline is the same for all pipelines:
</Tip>
* 1. Take an image and run it through a pre-conditioning processor.
* 2. Run the pre-processed image through the [`StableDiffusionControlNetPipeline`].
Let's have a look at a simple example using the [Canny Edge ControlNet](https://huggingface.co/lllyasviel/sd-controlnet-canny).
```python
from diffusers import StableDiffusionControlNetPipeline
from diffusers.utils import load_image
# Let's load the popular vermeer image
image = load_image(
"https://hf.co/datasets/huggingface/documentation-images/resolve/main/diffusers/input_image_vermeer.png"
)
```
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/input_image_vermeer.png)
Next, we process the image to get the canny image. This is step *1.* - running the pre-conditioning processor. The pre-conditioning processor is different for every ControlNet. Please see the model cards of the [official checkpoints](#controlnet-with-stable-diffusion-1.5) for more information about other models.
First, we need to install opencv:
```
pip install opencv-contrib-python
```
Next, let's also install all required Hugging Face libraries:
```
pip install diffusers transformers git+https://github.com/huggingface/accelerate.git
```
Then we can retrieve the canny edges of the image.
```python
import cv2
from PIL import Image
import numpy as np
image = np.array(image)
low_threshold = 100
high_threshold = 200
image = cv2.Canny(image, low_threshold, high_threshold)
image = image[:, :, None]
image = np.concatenate([image, image, image], axis=2)
canny_image = Image.fromarray(image)
```
Let's take a look at the processed image.
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/vermeer_canny_edged.png)
Now, we load the official [Stable Diffusion 1.5 Model](runwayml/stable-diffusion-v1-5) as well as the ControlNet for canny edges.
```py
from diffusers import StableDiffusionControlNetPipeline, ControlNetModel
import torch
controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16)
pipe = StableDiffusionControlNetPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
)
```
To speed-up things and reduce memory, let's enable model offloading and use the fast [`UniPCMultistepScheduler`].
```py
from diffusers import UniPCMultistepScheduler
pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
# this command loads the individual model components on GPU on-demand.
pipe.enable_model_cpu_offload()
```
Finally, we can run the pipeline:
```py
generator = torch.manual_seed(0)
out_image = pipe(
"disco dancer with colorful lights", num_inference_steps=20, generator=generator, image=canny_image
).images[0]
```
This should take only around 3-4 seconds on GPU (depending on hardware). The output image then looks as follows:
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/vermeer_disco_dancing.png)
**Note**: To see how to run all other ControlNet checkpoints, please have a look at [ControlNet with Stable Diffusion 1.5](#controlnet-with-stable-diffusion-1.5).
<!-- TODO: add space -->
## Combining multiple conditionings
Multiple ControlNet conditionings can be combined for a single image generation. Pass a list of ControlNets to the pipeline's constructor and a corresponding list of conditionings to `__call__`.
When combining conditionings, it is helpful to mask conditionings such that they do not overlap. In the example, we mask the middle of the canny map where the pose conditioning is located.
It can also be helpful to vary the `controlnet_conditioning_scales` to emphasize one conditioning over the other.
### Canny conditioning
The original image:
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/landscape.png"/>
Prepare the conditioning:
```python
from diffusers.utils import load_image
from PIL import Image
import cv2
import numpy as np
from diffusers.utils import load_image
canny_image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/landscape.png"
)
canny_image = np.array(canny_image)
low_threshold = 100
high_threshold = 200
canny_image = cv2.Canny(canny_image, low_threshold, high_threshold)
# zero out middle columns of image where pose will be overlayed
zero_start = canny_image.shape[1] // 4
zero_end = zero_start + canny_image.shape[1] // 2
canny_image[:, zero_start:zero_end] = 0
canny_image = canny_image[:, :, None]
canny_image = np.concatenate([canny_image, canny_image, canny_image], axis=2)
canny_image = Image.fromarray(canny_image)
```
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/controlnet/landscape_canny_masked.png"/>
### Openpose conditioning
The original image:
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/person.png" width=600/>
Prepare the conditioning:
```python
from controlnet_aux import OpenposeDetector
from diffusers.utils import load_image
openpose = OpenposeDetector.from_pretrained("lllyasviel/ControlNet")
openpose_image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/person.png"
)
openpose_image = openpose(openpose_image)
```
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/controlnet/person_pose.png" width=600/>
### Running ControlNet with multiple conditionings
```python
from diffusers import StableDiffusionControlNetPipeline, ControlNetModel, UniPCMultistepScheduler
import torch
controlnet = [
ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-openpose", torch_dtype=torch.float16),
ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny", torch_dtype=torch.float16),
]
pipe = StableDiffusionControlNetPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", controlnet=controlnet, torch_dtype=torch.float16
)
pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
pipe.enable_xformers_memory_efficient_attention()
pipe.enable_model_cpu_offload()
prompt = "a giant standing in a fantasy landscape, best quality"
negative_prompt = "monochrome, lowres, bad anatomy, worst quality, low quality"
generator = torch.Generator(device="cpu").manual_seed(1)
images = [openpose_image, canny_image]
image = pipe(
prompt,
images,
num_inference_steps=20,
generator=generator,
negative_prompt=negative_prompt,
controlnet_conditioning_scale=[1.0, 0.8],
).images[0]
image.save("./multi_controlnet_output.png")
```
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/controlnet/multi_controlnet_output.png" width=600/>
### Guess Mode
Guess Mode is [a ControlNet feature that was implemented](https://github.com/lllyasviel/ControlNet#guess-mode--non-prompt-mode) after the publication of [the paper](https://arxiv.org/abs/2302.05543). The description states:
>In this mode, the ControlNet encoder will try best to recognize the content of the input control map, like depth map, edge map, scribbles, etc, even if you remove all prompts.
#### The core implementation:
It adjusts the scale of the output residuals from ControlNet by a fixed ratio depending on the block depth. The shallowest DownBlock corresponds to `0.1`. As the blocks get deeper, the scale increases exponentially, and the scale for the output of the MidBlock becomes `1.0`.
Since the core implementation is just this, **it does not have any impact on prompt conditioning**. While it is common to use it without specifying any prompts, it is also possible to provide prompts if desired.
#### Usage:
Just specify `guess_mode=True` in the pipe() function. A `guidance_scale` between 3.0 and 5.0 is [recommended](https://github.com/lllyasviel/ControlNet#guess-mode--non-prompt-mode).
```py
from diffusers import StableDiffusionControlNetPipeline, ControlNetModel
import torch
controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-canny")
pipe = StableDiffusionControlNetPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", controlnet=controlnet).to(
"cuda"
)
image = pipe("", image=canny_image, guess_mode=True, guidance_scale=3.0).images[0]
image.save("guess_mode_generated.png")
```
#### Output image comparison:
Canny Control Example
|no guess_mode with prompt|guess_mode without prompt|
|---|---|
|<a href="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare_guess_mode/output_images/diffusers/output_bird_canny_0.png"><img width="128" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare_guess_mode/output_images/diffusers/output_bird_canny_0.png"/></a>|<a href="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare_guess_mode/output_images/diffusers/output_bird_canny_0_gm.png"><img width="128" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare_guess_mode/output_images/diffusers/output_bird_canny_0_gm.png"/></a>|
## Available checkpoints
ControlNet requires a *control image* in addition to the text-to-image *prompt*.
Each pretrained model is trained using a different conditioning method that requires different images for conditioning the generated outputs. For example, Canny edge conditioning requires the control image to be the output of a Canny filter, while depth conditioning requires the control image to be a depth map. See the overview and image examples below to know more.
All checkpoints can be found under the authors' namespace [lllyasviel](https://huggingface.co/lllyasviel).
**13.04.2024 Update**: The author has released improved controlnet checkpoints v1.1 - see [here](#controlnet-v1.1).
### ControlNet v1.0
| Model Name | Control Image Overview| Control Image Example | Generated Image Example |
|---|---|---|---|
|[lllyasviel/sd-controlnet-canny](https://huggingface.co/lllyasviel/sd-controlnet-canny)<br/> *Trained with canny edge detection* | A monochrome image with white edges on a black background.|<a href="https://huggingface.co/takuma104/controlnet_dev/blob/main/gen_compare/control_images/converted/control_bird_canny.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/control_images/converted/control_bird_canny.png"/></a>|<a href="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_bird_canny_1.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_bird_canny_1.png"/></a>|
|[lllyasviel/sd-controlnet-depth](https://huggingface.co/lllyasviel/sd-controlnet-depth)<br/> *Trained with Midas depth estimation* |A grayscale image with black representing deep areas and white representing shallow areas.|<a href="https://huggingface.co/takuma104/controlnet_dev/blob/main/gen_compare/control_images/converted/control_vermeer_depth.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/control_images/converted/control_vermeer_depth.png"/></a>|<a href="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_vermeer_depth_2.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_vermeer_depth_2.png"/></a>|
|[lllyasviel/sd-controlnet-hed](https://huggingface.co/lllyasviel/sd-controlnet-hed)<br/> *Trained with HED edge detection (soft edge)* |A monochrome image with white soft edges on a black background.|<a href="https://huggingface.co/takuma104/controlnet_dev/blob/main/gen_compare/control_images/converted/control_bird_hed.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/control_images/converted/control_bird_hed.png"/></a>|<a href="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_bird_hed_1.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_bird_hed_1.png"/></a> |
|[lllyasviel/sd-controlnet-mlsd](https://huggingface.co/lllyasviel/sd-controlnet-mlsd)<br/> *Trained with M-LSD line detection* |A monochrome image composed only of white straight lines on a black background.|<a href="https://huggingface.co/takuma104/controlnet_dev/blob/main/gen_compare/control_images/converted/control_room_mlsd.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/control_images/converted/control_room_mlsd.png"/></a>|<a href="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_room_mlsd_0.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_room_mlsd_0.png"/></a>|
|[lllyasviel/sd-controlnet-normal](https://huggingface.co/lllyasviel/sd-controlnet-normal)<br/> *Trained with normal map* |A [normal mapped](https://en.wikipedia.org/wiki/Normal_mapping) image.|<a href="https://huggingface.co/takuma104/controlnet_dev/blob/main/gen_compare/control_images/converted/control_human_normal.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/control_images/converted/control_human_normal.png"/></a>|<a href="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_human_normal_1.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_human_normal_1.png"/></a>|
|[lllyasviel/sd-controlnet-openpose](https://huggingface.co/lllyasviel/sd-controlnet_openpose)<br/> *Trained with OpenPose bone image* |A [OpenPose bone](https://github.com/CMU-Perceptual-Computing-Lab/openpose) image.|<a href="https://huggingface.co/takuma104/controlnet_dev/blob/main/gen_compare/control_images/converted/control_human_openpose.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/control_images/converted/control_human_openpose.png"/></a>|<a href="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_human_openpose_0.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_human_openpose_0.png"/></a>|
|[lllyasviel/sd-controlnet-scribble](https://huggingface.co/lllyasviel/sd-controlnet_scribble)<br/> *Trained with human scribbles* |A hand-drawn monochrome image with white outlines on a black background.|<a href="https://huggingface.co/takuma104/controlnet_dev/blob/main/gen_compare/control_images/converted/control_vermeer_scribble.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/control_images/converted/control_vermeer_scribble.png"/></a>|<a href="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_vermeer_scribble_0.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_vermeer_scribble_0.png"/></a> |
|[lllyasviel/sd-controlnet-seg](https://huggingface.co/lllyasviel/sd-controlnet_seg)<br/>*Trained with semantic segmentation* |An [ADE20K](https://groups.csail.mit.edu/vision/datasets/ADE20K/)'s segmentation protocol image.|<a href="https://huggingface.co/takuma104/controlnet_dev/blob/main/gen_compare/control_images/converted/control_room_seg.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/control_images/converted/control_room_seg.png"/></a>|<a href="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_room_seg_1.png"><img width="64" src="https://huggingface.co/takuma104/controlnet_dev/resolve/main/gen_compare/output_images/diffusers/output_room_seg_1.png"/></a> |
### ControlNet v1.1
| Model Name | Control Image Overview| Condition Image | Control Image Example | Generated Image Example |
|---|---|---|---|---|
|[lllyasviel/control_v11p_sd15_canny](https://huggingface.co/lllyasviel/control_v11p_sd15_canny)<br/> | *Trained with canny edge detection* | A monochrome image with white edges on a black background.|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_canny/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11p_sd15_canny/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_canny/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11p_sd15_canny/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11e_sd15_ip2p](https://huggingface.co/lllyasviel/control_v11e_sd15_ip2p)<br/> | *Trained with pixel to pixel instruction* | No condition .|<a href="https://huggingface.co/lllyasviel/control_v11e_sd15_ip2p/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11e_sd15_ip2p/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11e_sd15_ip2p/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11e_sd15_ip2p/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11p_sd15_inpaint](https://huggingface.co/lllyasviel/control_v11p_sd15_inpaint)<br/> | Trained with image inpainting | No condition.|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_inpaint/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11p_sd15_inpaint/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_inpaint/resolve/main/images/output.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11p_sd15_inpaint/resolve/main/images/output.png"/></a>|
|[lllyasviel/control_v11p_sd15_mlsd](https://huggingface.co/lllyasviel/control_v11p_sd15_mlsd)<br/> | Trained with multi-level line segment detection | An image with annotated line segments.|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_mlsd/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11p_sd15_mlsd/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_mlsd/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11p_sd15_mlsd/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11f1p_sd15_depth](https://huggingface.co/lllyasviel/control_v11f1p_sd15_depth)<br/> | Trained with depth estimation | An image with depth information, usually represented as a grayscale image.|<a href="https://huggingface.co/lllyasviel/control_v11f1p_sd15_depth/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11f1p_sd15_depth/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11f1p_sd15_depth/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11f1p_sd15_depth/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11p_sd15_normalbae](https://huggingface.co/lllyasviel/control_v11p_sd15_normalbae)<br/> | Trained with surface normal estimation | An image with surface normal information, usually represented as a color-coded image.|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_normalbae/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11p_sd15_normalbae/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_normalbae/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11p_sd15_normalbae/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11p_sd15_seg](https://huggingface.co/lllyasviel/control_v11p_sd15_seg)<br/> | Trained with image segmentation | An image with segmented regions, usually represented as a color-coded image.|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_seg/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11p_sd15_seg/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_seg/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11p_sd15_seg/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11p_sd15_lineart](https://huggingface.co/lllyasviel/control_v11p_sd15_lineart)<br/> | Trained with line art generation | An image with line art, usually black lines on a white background.|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_lineart/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11p_sd15_lineart/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_lineart/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11p_sd15_lineart/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11p_sd15s2_lineart_anime](https://huggingface.co/lllyasviel/control_v11p_sd15s2_lineart_anime)<br/> | Trained with anime line art generation | An image with anime-style line art.|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15s2_lineart_anime/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11p_sd15s2_lineart_anime/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15s2_lineart_anime/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11p_sd15s2_lineart_anime/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11p_sd15_openpose](https://huggingface.co/lllyasviel/control_v11p_sd15s2_lineart_anime)<br/> | Trained with human pose estimation | An image with human poses, usually represented as a set of keypoints or skeletons.|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_openpose/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11p_sd15_openpose/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_openpose/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11p_sd15_openpose/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11p_sd15_scribble](https://huggingface.co/lllyasviel/control_v11p_sd15_scribble)<br/> | Trained with scribble-based image generation | An image with scribbles, usually random or user-drawn strokes.|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_scribble/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11p_sd15_scribble/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_scribble/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11p_sd15_scribble/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11p_sd15_softedge](https://huggingface.co/lllyasviel/control_v11p_sd15_softedge)<br/> | Trained with soft edge image generation | An image with soft edges, usually to create a more painterly or artistic effect.|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_softedge/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11p_sd15_softedge/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11p_sd15_softedge/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11p_sd15_softedge/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11e_sd15_shuffle](https://huggingface.co/lllyasviel/control_v11e_sd15_shuffle)<br/> | Trained with image shuffling | An image with shuffled patches or regions.|<a href="https://huggingface.co/lllyasviel/control_v11e_sd15_shuffle/resolve/main/images/control.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11e_sd15_shuffle/resolve/main/images/control.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11e_sd15_shuffle/resolve/main/images/image_out.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11e_sd15_shuffle/resolve/main/images/image_out.png"/></a>|
|[lllyasviel/control_v11f1e_sd15_tile](https://huggingface.co/lllyasviel/control_v11f1e_sd15_tile)<br/> | Trained with image tiling | A blurry image or part of an image .|<a href="https://huggingface.co/lllyasviel/control_v11f1e_sd15_tile/resolve/main/images/original.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/lllyasviel/control_v11f1e_sd15_tile/resolve/main/images/original.png"/></a>|<a href="https://huggingface.co/lllyasviel/control_v11f1e_sd15_tile/resolve/main/images/output.png"><img width="64" src="https://huggingface.co/lllyasviel/control_v11f1e_sd15_tile/resolve/main/images/output.png"/></a>|
## StableDiffusionControlNetPipeline
[[autodoc]] StableDiffusionControlNetPipeline
@@ -66,15 +343,8 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
- disable_xformers_memory_efficient_attention
- load_textual_inversion
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
## FlaxStableDiffusionControlNetPipeline
[[autodoc]] FlaxStableDiffusionControlNetPipeline
- all
- __call__
## FlaxStableDiffusionControlNetPipelineOutput
[[autodoc]] pipelines.stable_diffusion.FlaxStableDiffusionPipelineOutput

View File

@@ -1,55 +0,0 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNet with Stable Diffusion XL
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
The abstract from the paper is:
*We present a neural network structure, ControlNet, to control pretrained large diffusion models to support additional input conditions. The ControlNet learns task-specific conditions in an end-to-end way, and the learning is robust even when the training dataset is small (< 50k). Moreover, training a ControlNet is as fast as fine-tuning a diffusion model, and the model can be trained on a personal devices. Alternatively, if powerful computation clusters are available, the model can scale to large amounts (millions to billions) of data. We report that large diffusion models like Stable Diffusion can be augmented with ControlNets to enable conditional inputs like edge maps, segmentation maps, keypoints, etc. This may enrich the methods to control large diffusion models and further facilitate related applications.*
You can find additional smaller Stable Diffusion XL (SDXL) ControlNet checkpoints from the 🤗 [Diffusers](https://huggingface.co/diffusers) Hub organization, and browse [community-trained](https://huggingface.co/models?other=stable-diffusion-xl&other=controlnet) checkpoints on the Hub.
<Tip warning={true}>
🧪 Many of the SDXL ControlNet checkpoints are experimental, and there is a lot of room for improvement. Feel free to open an [Issue](https://github.com/huggingface/diffusers/issues/new/choose) and leave us feedback on how we can improve!
</Tip>
If you don't see a checkpoint you're interested in, you can train your own SDXL ControlNet with our [training script](https://github.com/huggingface/diffusers/blob/main/examples/controlnet/README_sdxl.md).
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## StableDiffusionXLControlNetPipeline
[[autodoc]] StableDiffusionXLControlNetPipeline
- all
- __call__
## StableDiffusionXLControlNetImg2ImgPipeline
[[autodoc]] StableDiffusionXLControlNetImg2ImgPipeline
- all
- __call__
## StableDiffusionXLControlNetInpaintPipeline
[[autodoc]] StableDiffusionXLControlNetInpaintPipeline
- all
- __call__
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput

View File

@@ -20,7 +20,7 @@ The abstract from the paper is:
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -20,7 +20,7 @@ The original codebase of this implementation can be found at [Harmonai-org](http
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -22,7 +22,7 @@ The original codebase can be found at [hohonathanho/diffusion](https://github.co
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -24,32 +24,325 @@ This pipeline was contributed by [clarencechen](https://github.com/clarencechen)
## Tips
* The pipeline can generate masks that can be fed into other inpainting pipelines.
* In order to generate an image using this pipeline, both an image mask (source and target prompts can be manually specified or generated, and passed to [`~StableDiffusionDiffEditPipeline.generate_mask`])
and a set of partially inverted latents (generated using [`~StableDiffusionDiffEditPipeline.invert`]) _must_ be provided as arguments when calling the pipeline to generate the final edited image.
* The function [`~StableDiffusionDiffEditPipeline.generate_mask`] exposes two prompt arguments, `source_prompt` and `target_prompt`
* The pipeline can generate masks that can be fed into other inpainting pipelines. Check out the code examples below to know more.
* In order to generate an image using this pipeline, both an image mask (manually specified or generated using `generate_mask`)
and a set of partially inverted latents (generated using `invert`) _must_ be provided as arguments when calling the pipeline to generate the final edited image.
Refer to the code examples below for more details.
* The function `generate_mask` exposes two prompt arguments, `source_prompt` and `target_prompt`,
that let you control the locations of the semantic edits in the final image to be generated. Let's say,
you wanted to translate from "cat" to "dog". In this case, the edit direction will be "cat -> dog". To reflect
this in the generated mask, you simply have to set the embeddings related to the phrases including "cat" to
`source_prompt` and "dog" to `target_prompt`.
`source_prompt_embeds` and "dog" to `target_prompt_embeds`. Refer to the code example below for more details.
* When generating partially inverted latents using `invert`, assign a caption or text embedding describing the
overall image to the `prompt` argument to help guide the inverse latent sampling process. In most cases, the
source concept is sufficiently descriptive to yield good results, but feel free to explore alternatives.
source concept is sufficently descriptive to yield good results, but feel free to explore alternatives.
Please refer to [this code example](#generating-image-captions-for-inversion) for more details.
* When calling the pipeline to generate the final edited image, assign the source concept to `negative_prompt`
and the target concept to `prompt`. Taking the above example, you simply have to set the embeddings related to
the phrases including "cat" to `negative_prompt` and "dog" to `prompt`.
the phrases including "cat" to `negative_prompt_embeds` and "dog" to `prompt_embeds`. Refer to the code example
below for more details.
* If you wanted to reverse the direction in the example above, i.e., "dog -> cat", then it's recommended to:
* Swap the `source_prompt` and `target_prompt` in the arguments to `generate_mask`.
* Change the input prompt in [`~StableDiffusionDiffEditPipeline.invert`] to include "dog".
* Change the input prompt for `invert` to include "dog".
* Swap the `prompt` and `negative_prompt` in the arguments to call the pipeline to generate the final edited image.
* The source and target prompts, or their corresponding embeddings, can also be automatically generated. Please refer to the [DiffEdit](/using-diffusers/diffedit) guide for more details.
* Note that the source and target prompts, or their corresponding embeddings, can also be automatically generated. Please, refer to [this discussion](#generating-source-and-target-embeddings) for more details.
## Usage example
### Based on an input image with a caption
When the pipeline is conditioned on an input image, we first obtain partially inverted latents from the input image using a
`DDIMInverseScheduler` with the help of a caption. Then we generate an editing mask to identify relevant regions in the image using the source and target prompts. Finally,
the inverted noise and generated mask is used to start the generation process.
First, let's load our pipeline:
```py
import torch
from diffusers import DDIMScheduler, DDIMInverseScheduler, StableDiffusionDiffEditPipeline
sd_model_ckpt = "stabilityai/stable-diffusion-2-1"
pipeline = StableDiffusionDiffEditPipeline.from_pretrained(
sd_model_ckpt,
torch_dtype=torch.float16,
safety_checker=None,
)
pipeline.scheduler = DDIMScheduler.from_config(pipeline.scheduler.config)
pipeline.inverse_scheduler = DDIMInverseScheduler.from_config(pipeline.scheduler.config)
pipeline.enable_model_cpu_offload()
pipeline.enable_vae_slicing()
generator = torch.manual_seed(0)
```
Then, we load an input image to edit using our method:
```py
from diffusers.utils import load_image
img_url = "https://github.com/Xiang-cd/DiffEdit-stable-diffusion/raw/main/assets/origin.png"
raw_image = load_image(img_url).convert("RGB").resize((768, 768))
```
Then, we employ the source and target prompts to generate the editing mask:
```py
# See the "Generating source and target embeddings" section below to
# automate the generation of these captions with a pre-trained model like Flan-T5 as explained below.
source_prompt = "a bowl of fruits"
target_prompt = "a basket of fruits"
mask_image = pipeline.generate_mask(
image=raw_image,
source_prompt=source_prompt,
target_prompt=target_prompt,
generator=generator,
)
```
Then, we employ the caption and the input image to get the inverted latents:
```py
inv_latents = pipeline.invert(prompt=source_prompt, image=raw_image, generator=generator).latents
```
Now, generate the image with the inverted latents and semantically generated mask:
```py
image = pipeline(
prompt=target_prompt,
mask_image=mask_image,
image_latents=inv_latents,
generator=generator,
negative_prompt=source_prompt,
).images[0]
image.save("edited_image.png")
```
## Generating image captions for inversion
The authors originally used the source concept prompt as the caption for generating the partially inverted latents. However, we can also leverage open source and public image captioning models for the same purpose.
Below, we provide an end-to-end example with the [BLIP](https://huggingface.co/docs/transformers/model_doc/blip) model
for generating captions.
First, let's load our automatic image captioning model:
```py
import torch
from transformers import BlipForConditionalGeneration, BlipProcessor
captioner_id = "Salesforce/blip-image-captioning-base"
processor = BlipProcessor.from_pretrained(captioner_id)
model = BlipForConditionalGeneration.from_pretrained(captioner_id, torch_dtype=torch.float16, low_cpu_mem_usage=True)
```
Then, we define a utility to generate captions from an input image using the model:
```py
@torch.no_grad()
def generate_caption(images, caption_generator, caption_processor):
text = "a photograph of"
inputs = caption_processor(images, text, return_tensors="pt").to(device="cuda", dtype=caption_generator.dtype)
caption_generator.to("cuda")
outputs = caption_generator.generate(**inputs, max_new_tokens=128)
# offload caption generator
caption_generator.to("cpu")
caption = caption_processor.batch_decode(outputs, skip_special_tokens=True)[0]
return caption
```
Then, we load an input image for conditioning and obtain a suitable caption for it:
```py
from diffusers.utils import load_image
img_url = "https://github.com/Xiang-cd/DiffEdit-stable-diffusion/raw/main/assets/origin.png"
raw_image = load_image(img_url).convert("RGB").resize((768, 768))
caption = generate_caption(raw_image, model, processor)
```
Then, we employ the generated caption and the input image to get the inverted latents:
```py
from diffusers import DDIMInverseScheduler, DDIMScheduler
pipeline = StableDiffusionDiffEditPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-1", torch_dtype=torch.float16
)
pipeline = pipeline.to("cuda")
pipeline.enable_model_cpu_offload()
pipeline.enable_vae_slicing()
pipeline.scheduler = DDIMScheduler.from_config(pipeline.scheduler.config)
pipeline.inverse_scheduler = DDIMInverseScheduler.from_config(pipeline.scheduler.config)
generator = torch.manual_seed(0)
inv_latents = pipeline.invert(prompt=caption, image=raw_image, generator=generator).latents
```
Now, generate the image with the inverted latents and semantically generated mask from our source and target prompts:
```py
source_prompt = "a bowl of fruits"
target_prompt = "a basket of fruits"
mask_image = pipeline.generate_mask(
image=raw_image,
source_prompt=source_prompt,
target_prompt=target_prompt,
generator=generator,
)
image = pipeline(
prompt=target_prompt,
mask_image=mask_image,
image_latents=inv_latents,
generator=generator,
negative_prompt=source_prompt,
).images[0]
image.save("edited_image.png")
```
## Generating source and target embeddings
The authors originally required the user to manually provide the source and target prompts for discovering
edit directions. However, we can also leverage open source and public models for the same purpose.
Below, we provide an end-to-end example with the [Flan-T5](https://huggingface.co/docs/transformers/model_doc/flan-t5) model
for generating source an target embeddings.
**1. Load the generation model**:
```py
import torch
from transformers import AutoTokenizer, T5ForConditionalGeneration
tokenizer = AutoTokenizer.from_pretrained("google/flan-t5-xl")
model = T5ForConditionalGeneration.from_pretrained("google/flan-t5-xl", device_map="auto", torch_dtype=torch.float16)
```
**2. Construct a starting prompt**:
```py
source_concept = "bowl"
target_concept = "basket"
source_text = f"Provide a caption for images containing a {source_concept}. "
"The captions should be in English and should be no longer than 150 characters."
target_text = f"Provide a caption for images containing a {target_concept}. "
"The captions should be in English and should be no longer than 150 characters."
```
Here, we're interested in the "bowl -> basket" direction.
**3. Generate prompts**:
We can use a utility like so for this purpose.
```py
@torch.no_grad
def generate_prompts(input_prompt):
input_ids = tokenizer(input_prompt, return_tensors="pt").input_ids.to("cuda")
outputs = model.generate(
input_ids, temperature=0.8, num_return_sequences=16, do_sample=True, max_new_tokens=128, top_k=10
)
return tokenizer.batch_decode(outputs, skip_special_tokens=True)
```
And then we just call it to generate our prompts:
```py
source_prompts = generate_prompts(source_text)
target_prompts = generate_prompts(target_text)
```
We encourage you to play around with the different parameters supported by the
`generate()` method ([documentation](https://huggingface.co/docs/transformers/main/en/main_classes/text_generation#transformers.generation_tf_utils.TFGenerationMixin.generate)) for the generation quality you are looking for.
**4. Load the embedding model**:
Here, we need to use the same text encoder model used by the subsequent Stable Diffusion model.
```py
from diffusers import StableDiffusionDiffEditPipeline
pipeline = StableDiffusionDiffEditPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-1", torch_dtype=torch.float16
)
pipeline = pipeline.to("cuda")
pipeline.enable_model_cpu_offload()
pipeline.enable_vae_slicing()
generator = torch.manual_seed(0)
```
**5. Compute embeddings**:
```py
import torch
@torch.no_grad()
def embed_prompts(sentences, tokenizer, text_encoder, device="cuda"):
embeddings = []
for sent in sentences:
text_inputs = tokenizer(
sent,
padding="max_length",
max_length=tokenizer.model_max_length,
truncation=True,
return_tensors="pt",
)
text_input_ids = text_inputs.input_ids
prompt_embeds = text_encoder(text_input_ids.to(device), attention_mask=None)[0]
embeddings.append(prompt_embeds)
return torch.concatenate(embeddings, dim=0).mean(dim=0).unsqueeze(0)
source_embeddings = embed_prompts(source_prompts, pipeline.tokenizer, pipeline.text_encoder)
target_embeddings = embed_prompts(target_captions, pipeline.tokenizer, pipeline.text_encoder)
```
And you're done! Now, you can use these embeddings directly while calling the pipeline:
```py
from diffusers import DDIMInverseScheduler, DDIMScheduler
from diffusers.utils import load_image
pipeline.scheduler = DDIMScheduler.from_config(pipeline.scheduler.config)
pipeline.inverse_scheduler = DDIMInverseScheduler.from_config(pipeline.scheduler.config)
img_url = "https://github.com/Xiang-cd/DiffEdit-stable-diffusion/raw/main/assets/origin.png"
raw_image = load_image(img_url).convert("RGB").resize((768, 768))
mask_image = pipeline.generate_mask(
image=raw_image,
source_prompt_embeds=source_embeds,
target_prompt_embeds=target_embeds,
generator=generator,
)
inv_latents = pipeline.invert(
prompt_embeds=source_embeds,
image=raw_image,
generator=generator,
).latents
images = pipeline(
mask_image=mask_image,
image_latents=inv_latents,
prompt_embeds=target_embeddings,
negative_prompt_embeds=source_embeddings,
generator=generator,
).images
images[0].save("edited_image.png")
```
## StableDiffusionDiffEditPipeline
[[autodoc]] StableDiffusionDiffEditPipeline
- all
- generate_mask
- invert
- __call__
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
- __call__

View File

@@ -22,7 +22,7 @@ The original codebase can be found at [facebookresearch/dit](https://github.com/
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -7,60 +7,462 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Kandinsky 2.1
# Kandinsky
Kandinsky 2.1 is created by [Arseniy Shakhmatov](https://github.com/cene555), [Anton Razzhigaev](https://github.com/razzant), [Aleksandr Nikolich](https://github.com/AlexWortega), [Igor Pavlov](https://github.com/boomb0om), [Andrey Kuznetsov](https://github.com/kuznetsoffandrey) and [Denis Dimitrov](https://github.com/denndimitrov).
## Overview
The description from it's GitHub page is:
Kandinsky inherits best practices from [DALL-E 2](https://huggingface.co/papers/2204.06125) and [Latent Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/latent_diffusion), while introducing some new ideas.
*Kandinsky 2.1 inherits best practicies from Dall-E 2 and Latent diffusion, while introducing some new ideas. As text and image encoder it uses CLIP model and diffusion image prior (mapping) between latent spaces of CLIP modalities. This approach increases the visual performance of the model and unveils new horizons in blending images and text-guided image manipulation.*
It uses [CLIP](https://huggingface.co/docs/transformers/model_doc/clip) for encoding images and text, and a diffusion image prior (mapping) between latent spaces of CLIP modalities. This approach enhances the visual performance of the model and unveils new horizons in blending images and text-guided image manipulation.
The original codebase can be found at [ai-forever/Kandinsky-2](https://github.com/ai-forever/Kandinsky-2).
The Kandinsky model is created by [Arseniy Shakhmatov](https://github.com/cene555), [Anton Razzhigaev](https://github.com/razzant), [Aleksandr Nikolich](https://github.com/AlexWortega), [Igor Pavlov](https://github.com/boomb0om), [Andrey Kuznetsov](https://github.com/kuznetsoffandrey) and [Denis Dimitrov](https://github.com/denndimitrov). The original codebase can be found [here](https://github.com/ai-forever/Kandinsky-2)
<Tip>
Check out the [Kandinsky Community](https://huggingface.co/kandinsky-community) organization on the Hub for the official model checkpoints for tasks like text-to-image, image-to-image, and inpainting.
## Usage example
In the following, we will walk you through some examples of how to use the Kandinsky pipelines to create some visually aesthetic artwork.
### Text-to-Image Generation
For text-to-image generation, we need to use both [`KandinskyPriorPipeline`] and [`KandinskyPipeline`].
The first step is to encode text prompts with CLIP and then diffuse the CLIP text embeddings to CLIP image embeddings,
as first proposed in [DALL-E 2](https://cdn.openai.com/papers/dall-e-2.pdf).
Let's throw a fun prompt at Kandinsky to see what it comes up with.
```py
prompt = "A alien cheeseburger creature eating itself, claymation, cinematic, moody lighting"
```
First, let's instantiate the prior pipeline and the text-to-image pipeline. Both
pipelines are diffusion models.
```py
from diffusers import DiffusionPipeline
import torch
pipe_prior = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1-prior", torch_dtype=torch.float16)
pipe_prior.to("cuda")
t2i_pipe = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
t2i_pipe.to("cuda")
```
<Tip warning={true}>
By default, the text-to-image pipeline use [`DDIMScheduler`], you can change the scheduler to [`DDPMScheduler`]
```py
scheduler = DDPMScheduler.from_pretrained("kandinsky-community/kandinsky-2-1", subfolder="ddpm_scheduler")
t2i_pipe = DiffusionPipeline.from_pretrained(
"kandinsky-community/kandinsky-2-1", scheduler=scheduler, torch_dtype=torch.float16
)
t2i_pipe.to("cuda")
```
</Tip>
## KandinskyPriorPipeline
Now we pass the prompt through the prior to generate image embeddings. The prior
returns both the image embeddings corresponding to the prompt and negative/unconditional image
embeddings corresponding to an empty string.
```py
image_embeds, negative_image_embeds = pipe_prior(prompt, guidance_scale=1.0).to_tuple()
```
<Tip warning={true}>
The text-to-image pipeline expects both `image_embeds`, `negative_image_embeds` and the original
`prompt` as the text-to-image pipeline uses another text encoder to better guide the second diffusion
process of `t2i_pipe`.
By default, the prior returns unconditioned negative image embeddings corresponding to the negative prompt of `""`.
For better results, you can also pass a `negative_prompt` to the prior. This will increase the effective batch size
of the prior by a factor of 2.
```py
prompt = "A alien cheeseburger creature eating itself, claymation, cinematic, moody lighting"
negative_prompt = "low quality, bad quality"
image_embeds, negative_image_embeds = pipe_prior(prompt, negative_prompt, guidance_scale=1.0).to_tuple()
```
</Tip>
Next, we can pass the embeddings as well as the prompt to the text-to-image pipeline. Remember that
in case you are using a customized negative prompt, that you should pass this one also to the text-to-image pipelines
with `negative_prompt=negative_prompt`:
```py
image = t2i_pipe(
prompt, image_embeds=image_embeds, negative_image_embeds=negative_image_embeds, height=768, width=768
).images[0]
image.save("cheeseburger_monster.png")
```
One cheeseburger monster coming up! Enjoy!
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-docs/cheeseburger.png)
<Tip>
We also provide an end-to-end Kandinsky pipeline [`KandinskyCombinedPipeline`], which combines both the prior pipeline and text-to-image pipeline, and lets you perform inference in a single step. You can create the combined pipeline with the [`~AutoPipelineForTextToImage.from_pretrained`] method
```python
from diffusers import AutoPipelineForTextToImage
import torch
pipe = AutoPipelineForTextToImage.from_pretrained(
"kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16
)
pipe.enable_model_cpu_offload()
```
Under the hood, it will automatically load both [`KandinskyPriorPipeline`] and [`KandinskyPipeline`]. To generate images, you no longer need to call both pipelines and pass the outputs from one to another. You only need to call the combined pipeline once. You can set different `guidance_scale` and `num_inference_steps` for the prior pipeline with the `prior_guidance_scale` and `prior_num_inference_steps` arguments.
```python
prompt = "A alien cheeseburger creature eating itself, claymation, cinematic, moody lighting"
negative_prompt = "low quality, bad quality"
image = pipe(prompt=prompt, negative_prompt=negative_prompt, prior_guidance_scale =1.0, guidance_scacle = 4.0, height=768, width=768).images[0]
```
</Tip>
The Kandinsky model works extremely well with creative prompts. Here is some of the amazing art that can be created using the exact same process but with different prompts.
```python
prompt = "bird eye view shot of a full body woman with cyan light orange magenta makeup, digital art, long braided hair her face separated by makeup in the style of yin Yang surrealism, symmetrical face, real image, contrasting tone, pastel gradient background"
```
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-docs/hair.png)
```python
prompt = "A car exploding into colorful dust"
```
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-docs/dusts.png)
```python
prompt = "editorial photography of an organic, almost liquid smoke style armchair"
```
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-docs/smokechair.png)
```python
prompt = "birds eye view of a quilted paper style alien planet landscape, vibrant colours, Cinematic lighting"
```
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-docs/alienplanet.png)
### Text Guided Image-to-Image Generation
The same Kandinsky model weights can be used for text-guided image-to-image translation. In this case, just make sure to load the weights using the [`KandinskyImg2ImgPipeline`] pipeline.
**Note**: You can also directly move the weights of the text-to-image pipelines to the image-to-image pipelines
without loading them twice by making use of the [`~DiffusionPipeline.components`] function as explained [here](#converting-between-different-pipelines).
Let's download an image.
```python
from PIL import Image
import requests
from io import BytesIO
# download image
url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
response = requests.get(url)
original_image = Image.open(BytesIO(response.content)).convert("RGB")
original_image = original_image.resize((768, 512))
```
![img](https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg)
```python
import torch
from diffusers import KandinskyImg2ImgPipeline, KandinskyPriorPipeline
# create prior
pipe_prior = KandinskyPriorPipeline.from_pretrained(
"kandinsky-community/kandinsky-2-1-prior", torch_dtype=torch.float16
)
pipe_prior.to("cuda")
# create img2img pipeline
pipe = KandinskyImg2ImgPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
pipe.to("cuda")
prompt = "A fantasy landscape, Cinematic lighting"
negative_prompt = "low quality, bad quality"
image_embeds, negative_image_embeds = pipe_prior(prompt, negative_prompt).to_tuple()
out = pipe(
prompt,
image=original_image,
image_embeds=image_embeds,
negative_image_embeds=negative_image_embeds,
height=768,
width=768,
strength=0.3,
)
out.images[0].save("fantasy_land.png")
```
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-docs/img2img_fantasyland.png)
<Tip>
You can also use the [`KandinskyImg2ImgCombinedPipeline`] for end-to-end image-to-image generation with Kandinsky 2.1
```python
from diffusers import AutoPipelineForImage2Image
import torch
import requests
from io import BytesIO
from PIL import Image
import os
pipe = AutoPipelineForImage2Image.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
pipe.enable_model_cpu_offload()
prompt = "A fantasy landscape, Cinematic lighting"
negative_prompt = "low quality, bad quality"
url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
response = requests.get(url)
original_image = Image.open(BytesIO(response.content)).convert("RGB")
original_image.thumbnail((768, 768))
image = pipe(prompt=prompt, image=original_image, strength=0.3).images[0]
```
</Tip>
### Text Guided Inpainting Generation
You can use [`KandinskyInpaintPipeline`] to edit images. In this example, we will add a hat to the portrait of a cat.
```py
from diffusers import KandinskyInpaintPipeline, KandinskyPriorPipeline
from diffusers.utils import load_image
import torch
import numpy as np
pipe_prior = KandinskyPriorPipeline.from_pretrained(
"kandinsky-community/kandinsky-2-1-prior", torch_dtype=torch.float16
)
pipe_prior.to("cuda")
prompt = "a hat"
prior_output = pipe_prior(prompt)
pipe = KandinskyInpaintPipeline.from_pretrained("kandinsky-community/kandinsky-2-1-inpaint", torch_dtype=torch.float16)
pipe.to("cuda")
init_image = load_image(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main" "/kandinsky/cat.png"
)
mask = np.zeros((768, 768), dtype=np.float32)
# Let's mask out an area above the cat's head
mask[:250, 250:-250] = 1
out = pipe(
prompt,
image=init_image,
mask_image=mask,
**prior_output,
height=768,
width=768,
num_inference_steps=150,
)
image = out.images[0]
image.save("cat_with_hat.png")
```
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-docs/inpaint_cat_hat.png)
<Tip>
To use the [`KandinskyInpaintCombinedPipeline`] to perform end-to-end image inpainting generation, you can run below code instead
```python
from diffusers import AutoPipelineForInpainting
pipe = AutoPipelineForInpainting.from_pretrained("kandinsky-community/kandinsky-2-1-inpaint", torch_dtype=torch.float16)
pipe.enable_model_cpu_offload()
image = pipe(prompt=prompt, image=original_image, mask_image=mask).images[0]
```
</Tip>
🚨🚨🚨 __Breaking change for Kandinsky Mask Inpainting__ 🚨🚨🚨
We introduced a breaking change for Kandinsky inpainting pipeline in the following pull request: https://github.com/huggingface/diffusers/pull/4207. Previously we accepted a mask format where black pixels represent the masked-out area. This is inconsistent with all other pipelines in diffusers. We have changed the mask format in Knaindsky and now using white pixels instead.
Please upgrade your inpainting code to follow the above. If you are using Kandinsky Inpaint in production. You now need to change the mask to:
```python
# For PIL input
import PIL.ImageOps
mask = PIL.ImageOps.invert(mask)
# For PyTorch and Numpy input
mask = 1 - mask
```
### Interpolate
The [`KandinskyPriorPipeline`] also comes with a cool utility function that will allow you to interpolate the latent space of different images and texts super easily. Here is an example of how you can create an Impressionist-style portrait for your pet based on "The Starry Night".
Note that you can interpolate between texts and images - in the below example, we passed a text prompt "a cat" and two images to the `interplate` function, along with a `weights` variable containing the corresponding weights for each condition we interplate.
```python
from diffusers import KandinskyPriorPipeline, KandinskyPipeline
from diffusers.utils import load_image
import PIL
import torch
pipe_prior = KandinskyPriorPipeline.from_pretrained(
"kandinsky-community/kandinsky-2-1-prior", torch_dtype=torch.float16
)
pipe_prior.to("cuda")
img1 = load_image(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main" "/kandinsky/cat.png"
)
img2 = load_image(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main" "/kandinsky/starry_night.jpeg"
)
# add all the conditions we want to interpolate, can be either text or image
images_texts = ["a cat", img1, img2]
# specify the weights for each condition in images_texts
weights = [0.3, 0.3, 0.4]
# We can leave the prompt empty
prompt = ""
prior_out = pipe_prior.interpolate(images_texts, weights)
pipe = KandinskyPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
pipe.to("cuda")
image = pipe(prompt, **prior_out, height=768, width=768).images[0]
image.save("starry_cat.png")
```
![img](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/kandinsky-docs/starry_cat.png)
## Optimization
Running Kandinsky in inference requires running both a first prior pipeline: [`KandinskyPriorPipeline`]
and a second image decoding pipeline which is one of [`KandinskyPipeline`], [`KandinskyImg2ImgPipeline`], or [`KandinskyInpaintPipeline`].
The bulk of the computation time will always be the second image decoding pipeline, so when looking
into optimizing the model, one should look into the second image decoding pipeline.
When running with PyTorch < 2.0, we strongly recommend making use of [`xformers`](https://github.com/facebookresearch/xformers)
to speed-up the optimization. This can be done by simply running:
```py
from diffusers import DiffusionPipeline
import torch
t2i_pipe = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
t2i_pipe.enable_xformers_memory_efficient_attention()
```
When running on PyTorch >= 2.0, PyTorch's SDPA attention will automatically be used. For more information on
PyTorch's SDPA, feel free to have a look at [this blog post](https://pytorch.org/blog/accelerated-diffusers-pt-20/).
To have explicit control , you can also manually set the pipeline to use PyTorch's 2.0 efficient attention:
```py
from diffusers.models.attention_processor import AttnAddedKVProcessor2_0
t2i_pipe.unet.set_attn_processor(AttnAddedKVProcessor2_0())
```
The slowest and most memory intense attention processor is the default `AttnAddedKVProcessor` processor.
We do **not** recommend using it except for testing purposes or cases where very high determistic behaviour is desired.
You can set it with:
```py
from diffusers.models.attention_processor import AttnAddedKVProcessor
t2i_pipe.unet.set_attn_processor(AttnAddedKVProcessor())
```
With PyTorch >= 2.0, you can also use Kandinsky with `torch.compile` which depending
on your hardware can signficantly speed-up your inference time once the model is compiled.
To use Kandinsksy with `torch.compile`, you can do:
```py
t2i_pipe.unet.to(memory_format=torch.channels_last)
t2i_pipe.unet = torch.compile(t2i_pipe.unet, mode="reduce-overhead", fullgraph=True)
```
After compilation you should see a very fast inference time. For more information,
feel free to have a look at [Our PyTorch 2.0 benchmark](https://huggingface.co/docs/diffusers/main/en/optimization/torch2.0).
<Tip>
To generate images directly from a single pipeline, you can use [`KandinskyCombinedPipeline`], [`KandinskyImg2ImgCombinedPipeline`], [`KandinskyInpaintCombinedPipeline`].
These combined pipelines wrap the [`KandinskyPriorPipeline`] and [`KandinskyPipeline`], [`KandinskyImg2ImgPipeline`], [`KandinskyInpaintPipeline`] respectively into a single
pipeline for a simpler user experience
</Tip>
## Available Pipelines:
| Pipeline | Tasks |
|---|---|
| [pipeline_kandinsky.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky/pipeline_kandinsky.py) | *Text-to-Image Generation* |
| [pipeline_kandinsky_combined.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky_combined.py) | *End-to-end Text-to-Image, image-to-image, Inpainting Generation* |
| [pipeline_kandinsky_inpaint.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky/pipeline_kandinsky_inpaint.py) | *Image-Guided Image Generation* |
| [pipeline_kandinsky_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky/pipeline_kandinsky_img2img.py) | *Image-Guided Image Generation* |
### KandinskyPriorPipeline
[[autodoc]] KandinskyPriorPipeline
- all
- __call__
- interpolate
## KandinskyPipeline
### KandinskyPipeline
[[autodoc]] KandinskyPipeline
- all
- __call__
## KandinskyCombinedPipeline
[[autodoc]] KandinskyCombinedPipeline
- all
- __call__
## KandinskyImg2ImgPipeline
### KandinskyImg2ImgPipeline
[[autodoc]] KandinskyImg2ImgPipeline
- all
- __call__
## KandinskyImg2ImgCombinedPipeline
[[autodoc]] KandinskyImg2ImgCombinedPipeline
- all
- __call__
## KandinskyInpaintPipeline
### KandinskyInpaintPipeline
[[autodoc]] KandinskyInpaintPipeline
- all
- __call__
## KandinskyInpaintCombinedPipeline
### KandinskyCombinedPipeline
[[autodoc]] KandinskyCombinedPipeline
- all
- __call__
### KandinskyImg2ImgCombinedPipeline
[[autodoc]] KandinskyImg2ImgCombinedPipeline
- all
- __call__
### KandinskyInpaintCombinedPipeline
[[autodoc]] KandinskyInpaintCombinedPipeline
- all

View File

@@ -9,77 +9,348 @@ specific language governing permissions and limitations under the License.
# Kandinsky 2.2
Kandinsky 2.1 is created by [Arseniy Shakhmatov](https://github.com/cene555), [Anton Razzhigaev](https://github.com/razzant), [Aleksandr Nikolich](https://github.com/AlexWortega), [Igor Pavlov](https://github.com/boomb0om), [Andrey Kuznetsov](https://github.com/kuznetsoffandrey) and [Denis Dimitrov](https://github.com/denndimitrov).
The Kandinsky 2.2 release includes robust new text-to-image models that support text-to-image generation, image-to-image generation, image interpolation, and text-guided image inpainting. The general workflow to perform these tasks using Kandinsky 2.2 is the same as in Kandinsky 2.1. First, you will need to use a prior pipeline to generate image embeddings based on your text prompt, and then use one of the image decoding pipelines to generate the output image. The only difference is that in Kandinsky 2.2, all of the decoding pipelines no longer accept the `prompt` input, and the image generation process is conditioned with only `image_embeds` and `negative_image_embeds`.
The description from it's GitHub page is:
Same as with Kandinsky 2.1, the easiest way to perform text-to-image generation is to use the combined Kandinsky pipeline. This process is exactly the same as Kandinsky 2.1. All you need to do is to replace the Kandinsky 2.1 checkpoint with 2.2.
*Kandinsky 2.2 brings substantial improvements upon its predecessor, Kandinsky 2.1, by introducing a new, more powerful image encoder - CLIP-ViT-G and the ControlNet support. The switch to CLIP-ViT-G as the image encoder significantly increases the model's capability to generate more aesthetic pictures and better understand text, thus enhancing the model's overall performance. The addition of the ControlNet mechanism allows the model to effectively control the process of generating images. This leads to more accurate and visually appealing outputs and opens new possibilities for text-guided image manipulation.*
```python
from diffusers import AutoPipelineForText2Image
import torch
The original codebase can be found at [ai-forever/Kandinsky-2](https://github.com/ai-forever/Kandinsky-2).
pipe = AutoPipelineForText2Image.from_pretrained("kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16)
pipe.enable_model_cpu_offload()
prompt = "A alien cheeseburger creature eating itself, claymation, cinematic, moody lighting"
negative_prompt = "low quality, bad quality"
image = pipe(prompt=prompt, negative_prompt=negative_prompt, prior_guidance_scale =1.0, height=768, width=768).images[0]
```
Now, let's look at an example where we take separate steps to run the prior pipeline and text-to-image pipeline. This way, we can understand what's happening under the hood and how Kandinsky 2.2 differs from Kandinsky 2.1.
First, let's create the prior pipeline and text-to-image pipeline with Kandinsky 2.2 checkpoints.
```python
from diffusers import DiffusionPipeline
import torch
pipe_prior = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-2-prior", torch_dtype=torch.float16)
pipe_prior.to("cuda")
t2i_pipe = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-2-decoder", torch_dtype=torch.float16)
t2i_pipe.to("cuda")
```
You can then use `pipe_prior` to generate image embeddings.
```python
prompt = "portrait of a women, blue eyes, cinematic"
negative_prompt = "low quality, bad quality"
image_embeds, negative_image_embeds = pipe_prior(prompt, guidance_scale=1.0).to_tuple()
```
Now you can pass these embeddings to the text-to-image pipeline. When using Kandinsky 2.2 you don't need to pass the `prompt` (but you do with the previous version, Kandinsky 2.1).
```
image = t2i_pipe(image_embeds=image_embeds, negative_image_embeds=negative_image_embeds, height=768, width=768).images[
0
]
image.save("portrait.png")
```
![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/kandinskyv22/%20blue%20eyes.png)
We used the text-to-image pipeline as an example, but the same process applies to all decoding pipelines in Kandinsky 2.2. For more information, please refer to our API section for each pipeline.
### Text-to-Image Generation with ControlNet Conditioning
In the following, we give a simple example of how to use [`KandinskyV22ControlnetPipeline`] to add control to the text-to-image generation with a depth image.
First, let's take an image and extract its depth map.
```python
from diffusers.utils import load_image
img = load_image(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/kandinskyv22/cat.png"
).resize((768, 768))
```
![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/kandinskyv22/cat.png)
We can use the `depth-estimation` pipeline from transformers to process the image and retrieve its depth map.
```python
import torch
import numpy as np
from transformers import pipeline
from diffusers.utils import load_image
def make_hint(image, depth_estimator):
image = depth_estimator(image)["depth"]
image = np.array(image)
image = image[:, :, None]
image = np.concatenate([image, image, image], axis=2)
detected_map = torch.from_numpy(image).float() / 255.0
hint = detected_map.permute(2, 0, 1)
return hint
depth_estimator = pipeline("depth-estimation")
hint = make_hint(img, depth_estimator).unsqueeze(0).half().to("cuda")
```
Now, we load the prior pipeline and the text-to-image controlnet pipeline
```python
from diffusers import KandinskyV22PriorPipeline, KandinskyV22ControlnetPipeline
pipe_prior = KandinskyV22PriorPipeline.from_pretrained(
"kandinsky-community/kandinsky-2-2-prior", torch_dtype=torch.float16
)
pipe_prior = pipe_prior.to("cuda")
pipe = KandinskyV22ControlnetPipeline.from_pretrained(
"kandinsky-community/kandinsky-2-2-controlnet-depth", torch_dtype=torch.float16
)
pipe = pipe.to("cuda")
```
We pass the prompt and negative prompt through the prior to generate image embeddings
```python
prompt = "A robot, 4k photo"
negative_prior_prompt = "lowres, text, error, cropped, worst quality, low quality, jpeg artifacts, ugly, duplicate, morbid, mutilated, out of frame, extra fingers, mutated hands, poorly drawn hands, poorly drawn face, mutation, deformed, blurry, dehydrated, bad anatomy, bad proportions, extra limbs, cloned face, disfigured, gross proportions, malformed limbs, missing arms, missing legs, extra arms, extra legs, fused fingers, too many fingers, long neck, username, watermark, signature"
generator = torch.Generator(device="cuda").manual_seed(43)
image_emb, zero_image_emb = pipe_prior(
prompt=prompt, negative_prompt=negative_prior_prompt, generator=generator
).to_tuple()
```
Now we can pass the image embeddings and the depth image we extracted to the controlnet pipeline. With Kandinsky 2.2, only prior pipelines accept `prompt` input. You do not need to pass the prompt to the controlnet pipeline.
```python
images = pipe(
image_embeds=image_emb,
negative_image_embeds=zero_image_emb,
hint=hint,
num_inference_steps=50,
generator=generator,
height=768,
width=768,
).images
images[0].save("robot_cat.png")
```
The output image looks as follow:
![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/kandinskyv22/robot_cat_text2img.png)
### Image-to-Image Generation with ControlNet Conditioning
Kandinsky 2.2 also includes a [`KandinskyV22ControlnetImg2ImgPipeline`] that will allow you to add control to the image generation process with both the image and its depth map. This pipeline works really well with [`KandinskyV22PriorEmb2EmbPipeline`], which generates image embeddings based on both a text prompt and an image.
For our robot cat example, we will pass the prompt and cat image together to the prior pipeline to generate an image embedding. We will then use that image embedding and the depth map of the cat to further control the image generation process.
We can use the same cat image and its depth map from the last example.
```python
import torch
import numpy as np
from diffusers import KandinskyV22PriorEmb2EmbPipeline, KandinskyV22ControlnetImg2ImgPipeline
from diffusers.utils import load_image
from transformers import pipeline
img = load_image(
"https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main" "/kandinskyv22/cat.png"
).resize((768, 768))
def make_hint(image, depth_estimator):
image = depth_estimator(image)["depth"]
image = np.array(image)
image = image[:, :, None]
image = np.concatenate([image, image, image], axis=2)
detected_map = torch.from_numpy(image).float() / 255.0
hint = detected_map.permute(2, 0, 1)
return hint
depth_estimator = pipeline("depth-estimation")
hint = make_hint(img, depth_estimator).unsqueeze(0).half().to("cuda")
pipe_prior = KandinskyV22PriorEmb2EmbPipeline.from_pretrained(
"kandinsky-community/kandinsky-2-2-prior", torch_dtype=torch.float16
)
pipe_prior = pipe_prior.to("cuda")
pipe = KandinskyV22ControlnetImg2ImgPipeline.from_pretrained(
"kandinsky-community/kandinsky-2-2-controlnet-depth", torch_dtype=torch.float16
)
pipe = pipe.to("cuda")
prompt = "A robot, 4k photo"
negative_prior_prompt = "lowres, text, error, cropped, worst quality, low quality, jpeg artifacts, ugly, duplicate, morbid, mutilated, out of frame, extra fingers, mutated hands, poorly drawn hands, poorly drawn face, mutation, deformed, blurry, dehydrated, bad anatomy, bad proportions, extra limbs, cloned face, disfigured, gross proportions, malformed limbs, missing arms, missing legs, extra arms, extra legs, fused fingers, too many fingers, long neck, username, watermark, signature"
generator = torch.Generator(device="cuda").manual_seed(43)
# run prior pipeline
img_emb = pipe_prior(prompt=prompt, image=img, strength=0.85, generator=generator)
negative_emb = pipe_prior(prompt=negative_prior_prompt, image=img, strength=1, generator=generator)
# run controlnet img2img pipeline
images = pipe(
image=img,
strength=0.5,
image_embeds=img_emb.image_embeds,
negative_image_embeds=negative_emb.image_embeds,
hint=hint,
num_inference_steps=50,
generator=generator,
height=768,
width=768,
).images
images[0].save("robot_cat.png")
```
Here is the output. Compared with the output from our text-to-image controlnet example, it kept a lot more cat facial details from the original image and worked into the robot style we asked for.
![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/kandinskyv22/robot_cat.png)
## Optimization
Running Kandinsky in inference requires running both a first prior pipeline: [`KandinskyPriorPipeline`]
and a second image decoding pipeline which is one of [`KandinskyPipeline`], [`KandinskyImg2ImgPipeline`], or [`KandinskyInpaintPipeline`].
The bulk of the computation time will always be the second image decoding pipeline, so when looking
into optimizing the model, one should look into the second image decoding pipeline.
When running with PyTorch < 2.0, we strongly recommend making use of [`xformers`](https://github.com/facebookresearch/xformers)
to speed-up the optimization. This can be done by simply running:
```py
from diffusers import DiffusionPipeline
import torch
t2i_pipe = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
t2i_pipe.enable_xformers_memory_efficient_attention()
```
When running on PyTorch >= 2.0, PyTorch's SDPA attention will automatically be used. For more information on
PyTorch's SDPA, feel free to have a look at [this blog post](https://pytorch.org/blog/accelerated-diffusers-pt-20/).
To have explicit control , you can also manually set the pipeline to use PyTorch's 2.0 efficient attention:
```py
from diffusers.models.attention_processor import AttnAddedKVProcessor2_0
t2i_pipe.unet.set_attn_processor(AttnAddedKVProcessor2_0())
```
The slowest and most memory intense attention processor is the default `AttnAddedKVProcessor` processor.
We do **not** recommend using it except for testing purposes or cases where very high determistic behaviour is desired.
You can set it with:
```py
from diffusers.models.attention_processor import AttnAddedKVProcessor
t2i_pipe.unet.set_attn_processor(AttnAddedKVProcessor())
```
With PyTorch >= 2.0, you can also use Kandinsky with `torch.compile` which depending
on your hardware can signficantly speed-up your inference time once the model is compiled.
To use Kandinsksy with `torch.compile`, you can do:
```py
t2i_pipe.unet.to(memory_format=torch.channels_last)
t2i_pipe.unet = torch.compile(t2i_pipe.unet, mode="reduce-overhead", fullgraph=True)
```
After compilation you should see a very fast inference time. For more information,
feel free to have a look at [Our PyTorch 2.0 benchmark](https://huggingface.co/docs/diffusers/main/en/optimization/torch2.0).
<Tip>
Check out the [Kandinsky Community](https://huggingface.co/kandinsky-community) organization on the Hub for the official model checkpoints for tasks like text-to-image, image-to-image, and inpainting.
To generate images directly from a single pipeline, you can use [`KandinskyV22CombinedPipeline`], [`KandinskyV22Img2ImgCombinedPipeline`], [`KandinskyV22InpaintCombinedPipeline`].
These combined pipelines wrap the [`KandinskyV22PriorPipeline`] and [`KandinskyV22Pipeline`], [`KandinskyV22Img2ImgPipeline`], [`KandinskyV22InpaintPipeline`] respectively into a single
pipeline for a simpler user experience
</Tip>
## KandinskyV22PriorPipeline
## Available Pipelines:
| Pipeline | Tasks |
|---|---|
| [pipeline_kandinsky2_2.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2.py) | *Text-to-Image Generation* |
| [pipeline_kandinsky2_2_combined.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2_combined.py) | *End-to-end Text-to-Image, image-to-image, Inpainting Generation* |
| [pipeline_kandinsky2_2_inpaint.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2_inpaint.py) | *Image-Guided Image Generation* |
| [pipeline_kandinsky2_2_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2_img2img.py) | *Image-Guided Image Generation* |
| [pipeline_kandinsky2_2_controlnet.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2_controlnet.py) | *Image-Guided Image Generation* |
| [pipeline_kandinsky2_2_controlnet_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/kandinsky2_2/pipeline_kandinsky2_2_controlnet_img2img.py) | *Image-Guided Image Generation* |
### KandinskyV22Pipeline
[[autodoc]] KandinskyV22Pipeline
- all
- __call__
### KandinskyV22ControlnetPipeline
[[autodoc]] KandinskyV22ControlnetPipeline
- all
- __call__
### KandinskyV22ControlnetImg2ImgPipeline
[[autodoc]] KandinskyV22ControlnetImg2ImgPipeline
- all
- __call__
### KandinskyV22Img2ImgPipeline
[[autodoc]] KandinskyV22Img2ImgPipeline
- all
- __call__
### KandinskyV22InpaintPipeline
[[autodoc]] KandinskyV22InpaintPipeline
- all
- __call__
### KandinskyV22PriorPipeline
[[autodoc]] KandinskyV22PriorPipeline
- all
- __call__
- interpolate
## KandinskyV22Pipeline
[[autodoc]] KandinskyV22Pipeline
- all
- __call__
## KandinskyV22CombinedPipeline
[[autodoc]] KandinskyV22CombinedPipeline
- all
- __call__
## KandinskyV22ControlnetPipeline
[[autodoc]] KandinskyV22ControlnetPipeline
- all
- __call__
## KandinskyV22PriorEmb2EmbPipeline
### KandinskyV22PriorEmb2EmbPipeline
[[autodoc]] KandinskyV22PriorEmb2EmbPipeline
- all
- __call__
- interpolate
## KandinskyV22Img2ImgPipeline
### KandinskyV22CombinedPipeline
[[autodoc]] KandinskyV22Img2ImgPipeline
[[autodoc]] KandinskyV22CombinedPipeline
- all
- __call__
## KandinskyV22Img2ImgCombinedPipeline
### KandinskyV22Img2ImgCombinedPipeline
[[autodoc]] KandinskyV22Img2ImgCombinedPipeline
- all
- __call__
## KandinskyV22ControlnetImg2ImgPipeline
[[autodoc]] KandinskyV22ControlnetImg2ImgPipeline
- all
- __call__
## KandinskyV22InpaintPipeline
[[autodoc]] KandinskyV22InpaintPipeline
- all
- __call__
## KandinskyV22InpaintCombinedPipeline
### KandinskyV22InpaintCombinedPipeline
[[autodoc]] KandinskyV22InpaintCombinedPipeline
- all

View File

@@ -1,40 +0,0 @@
# Latent Consistency Models
Latent Consistency Models (LCMs) were proposed in [Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference](https://arxiv.org/abs/2310.04378) by Simian Luo, Yiqin Tan, Longbo Huang, Jian Li, and Hang Zhao.
The abstract of the [paper](https://arxiv.org/pdf/2310.04378.pdf) is as follows:
*Latent Diffusion models (LDMs) have achieved remarkable results in synthesizing high-resolution images. However, the iterative sampling process is computationally intensive and leads to slow generation. Inspired by Consistency Models (song et al.), we propose Latent Consistency Models (LCMs), enabling swift inference with minimal steps on any pre-trained LDMs, including Stable Diffusion (rombach et al). Viewing the guided reverse diffusion process as solving an augmented probability flow ODE (PF-ODE), LCMs are designed to directly predict the solution of such ODE in latent space, mitigating the need for numerous iterations and allowing rapid, high-fidelity sampling. Efficiently distilled from pre-trained classifier-free guided diffusion models, a high-quality 768 x 768 2~4-step LCM takes only 32 A100 GPU hours for training. Furthermore, we introduce Latent Consistency Fine-tuning (LCF), a novel method that is tailored for fine-tuning LCMs on customized image datasets. Evaluation on the LAION-5B-Aesthetics dataset demonstrates that LCMs achieve state-of-the-art text-to-image generation performance with few-step inference.*
A demo for the [SimianLuo/LCM_Dreamshaper_v7](https://huggingface.co/SimianLuo/LCM_Dreamshaper_v7) checkpoint can be found [here](https://huggingface.co/spaces/SimianLuo/Latent_Consistency_Model).
The pipelines were contributed by [luosiallen](https://luosiallen.github.io/), [nagolinc](https://github.com/nagolinc), and [dg845](https://github.com/dg845).
## LatentConsistencyModelPipeline
[[autodoc]] LatentConsistencyModelPipeline
- all
- __call__
- enable_freeu
- disable_freeu
- enable_vae_slicing
- disable_vae_slicing
- enable_vae_tiling
- disable_vae_tiling
## LatentConsistencyModelImg2ImgPipeline
[[autodoc]] LatentConsistencyModelImg2ImgPipeline
- all
- __call__
- enable_freeu
- disable_freeu
- enable_vae_slicing
- disable_vae_slicing
- enable_vae_tiling
- disable_vae_tiling
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput

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@@ -22,7 +22,7 @@ The original codebase can be found at [Compvis/latent-diffusion](https://github.
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -22,7 +22,7 @@ The original codebase can be found at [CompVis/latent-diffusion](https://github.
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

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@@ -22,7 +22,7 @@ You can find additional information about model editing on the [project page](ht
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

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@@ -1,55 +0,0 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# MusicLDM
MusicLDM was proposed in [MusicLDM: Enhancing Novelty in Text-to-Music Generation Using Beat-Synchronous Mixup Strategies](https://huggingface.co/papers/2308.01546) by Ke Chen, Yusong Wu, Haohe Liu, Marianna Nezhurina, Taylor Berg-Kirkpatrick, Shlomo Dubnov.
MusicLDM takes a text prompt as input and predicts the corresponding music sample.
Inspired by [Stable Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/overview) and [AudioLDM](https://huggingface.co/docs/diffusers/api/pipelines/audioldm/overview),
MusicLDM is a text-to-music _latent diffusion model (LDM)_ that learns continuous audio representations from [CLAP](https://huggingface.co/docs/transformers/main/model_doc/clap)
latents.
MusicLDM is trained on a corpus of 466 hours of music data. Beat-synchronous data augmentation strategies are applied to
the music samples, both in the time domain and in the latent space. Using beat-synchronous data augmentation strategies
encourages the model to interpolate between the training samples, but stay within the domain of the training data. The
result is generated music that is more diverse while staying faithful to the corresponding style.
The abstract of the paper is the following:
*In this paper, we present MusicLDM, a state-of-the-art text-to-music model that adapts Stable Diffusion and AudioLDM architectures to the music domain. We achieve this by retraining the contrastive language-audio pretraining model (CLAP) and the Hifi-GAN vocoder, as components of MusicLDM, on a collection of music data samples. Then, we leverage a beat tracking model and propose two different mixup strategies for data augmentation: beat-synchronous audio mixup and beat-synchronous latent mixup, to encourage the model to generate music more diverse while still staying faithful to the corresponding style.*
This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi).
## Tips
When constructing a prompt, keep in mind:
* Descriptive prompt inputs work best; use adjectives to describe the sound (for example, "high quality" or "clear") and make the prompt context specific where possible (e.g. "melodic techno with a fast beat and synths" works better than "techno").
* Using a *negative prompt* can significantly improve the quality of the generated audio. Try using a negative prompt of "low quality, average quality".
During inference:
* The _quality_ of the generated audio sample can be controlled by the `num_inference_steps` argument; higher steps give higher quality audio at the expense of slower inference.
* Multiple waveforms can be generated in one go: set `num_waveforms_per_prompt` to a value greater than 1 to enable. Automatic scoring will be performed between the generated waveforms and prompt text, and the audios ranked from best to worst accordingly.
* The _length_ of the generated audio sample can be controlled by varying the `audio_length_in_s` argument.
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## MusicLDMPipeline
[[autodoc]] MusicLDMPipeline
- all
- __call__

View File

@@ -12,74 +12,16 @@ specific language governing permissions and limitations under the License.
# Pipelines
Pipelines provide a simple way to run state-of-the-art diffusion models in inference by bundling all of the necessary components (multiple independently-trained models, schedulers, and processors) into a single end-to-end class. Pipelines are flexible and they can be adapted to use different schedulers or even model components.
Pipelines provide a simple way to run state-of-the-art diffusion models in inference by bundling all of the necessary components (multiple independently-trained models, schedulers, and processors) into a single end-to-end class. Pipelines are flexible and they can be adapted to use different scheduler or even model components.
All pipelines are built from the base [`DiffusionPipeline`] class which provides basic functionality for loading, downloading, and saving all the components. Specific pipeline types (for example [`StableDiffusionPipeline`]) loaded with [`~DiffusionPipeline.from_pretrained`] are automatically detected and the pipeline components are loaded and passed to the `__init__` function of the pipeline.
All pipelines are built from the base [`DiffusionPipeline`] class which provides basic functionality for loading, downloading, and saving all the components.
<Tip warning={true}>
You shouldn't use the [`DiffusionPipeline`] class for training. Individual components (for example, [`UNet2DModel`] and [`UNet2DConditionModel`]) of diffusion pipelines are usually trained individually, so we suggest directly working with them instead.
<br>
Pipelines do not offer any training functionality. You'll notice PyTorch's autograd is disabled by decorating the [`~DiffusionPipeline.__call__`] method with a [`torch.no_grad`](https://pytorch.org/docs/stable/generated/torch.no_grad.html) decorator because pipelines should not be used for training. If you're interested in training, please take a look at the [Training](../../training/overview) guides instead!
Pipelines do not offer any training functionality. You'll notice PyTorch's autograd is disabled by decorating the [`~DiffusionPipeline.__call__`] method with a [`torch.no_grad`](https://pytorch.org/docs/stable/generated/torch.no_grad.html) decorator because pipelines should not be used for training. If you're interested in training, please take a look at the [Training](../traininig/overview) guides instead!
</Tip>
The table below lists all the pipelines currently available in 🤗 Diffusers and the tasks they support. Click on a pipeline to view its abstract and published paper.
| Pipeline | Tasks |
|---|---|
| [AltDiffusion](alt_diffusion) | image2image |
| [Attend-and-Excite](attend_and_excite) | text2image |
| [Audio Diffusion](audio_diffusion) | image2audio |
| [AudioLDM](audioldm) | text2audio |
| [AudioLDM2](audioldm2) | text2audio |
| [BLIP Diffusion](blip_diffusion) | text2image |
| [Consistency Models](consistency_models) | unconditional image generation |
| [ControlNet](controlnet) | text2image, image2image, inpainting |
| [ControlNet with Stable Diffusion XL](controlnet_sdxl) | text2image |
| [Cycle Diffusion](cycle_diffusion) | image2image |
| [Dance Diffusion](dance_diffusion) | unconditional audio generation |
| [DDIM](ddim) | unconditional image generation |
| [DDPM](ddpm) | unconditional image generation |
| [DeepFloyd IF](deepfloyd_if) | text2image, image2image, inpainting, super-resolution |
| [DiffEdit](diffedit) | inpainting |
| [DiT](dit) | text2image |
| [GLIGEN](gligen) | text2image |
| [InstructPix2Pix](pix2pix) | image editing |
| [Kandinsky](kandinsky) | text2image, image2image, inpainting, interpolation |
| [Kandinsky 2.2](kandinsky_v22) | text2image, image2image, inpainting |
| [Latent Diffusion](latent_diffusion) | text2image, super-resolution |
| [LDM3D](ldm3d_diffusion) | text2image, text-to-3D |
| [MultiDiffusion](panorama) | text2image |
| [MusicLDM](musicldm) | text2audio |
| [PaintByExample](paint_by_example) | inpainting |
| [ParaDiGMS](paradigms) | text2image |
| [Pix2Pix Zero](pix2pix_zero) | image editing |
| [PNDM](pndm) | unconditional image generation |
| [RePaint](repaint) | inpainting |
| [ScoreSdeVe](score_sde_ve) | unconditional image generation |
| [Self-Attention Guidance](self_attention_guidance) | text2image |
| [Semantic Guidance](semantic_stable_diffusion) | text2image |
| [Shap-E](shap_e) | text-to-3D, image-to-3D |
| [Spectrogram Diffusion](spectrogram_diffusion) | |
| [Stable Diffusion](stable_diffusion/overview) | text2image, image2image, depth2image, inpainting, image variation, latent upscaler, super-resolution |
| [Stable Diffusion Model Editing](model_editing) | model editing |
| [Stable Diffusion XL](stable_diffusion_xl) | text2image, image2image, inpainting |
| [Stable unCLIP](stable_unclip) | text2image, image variation |
| [KarrasVe](karras_ve) | unconditional image generation |
| [T2I Adapter](adapter) | text2image |
| [Text2Video](text_to_video) | text2video, video2video |
| [Text2Video Zero](text_to_video_zero) | text2video |
| [UnCLIP](unclip) | text2image, image variation |
| [Unconditional Latent Diffusion](latent_diffusion_uncond) | unconditional image generation |
| [UniDiffuser](unidiffuser) | text2image, image2text, image variation, text variation, unconditional image generation, unconditional audio generation |
| [Value-guided planning](value_guided_sampling) | value guided sampling |
| [Versatile Diffusion](versatile_diffusion) | text2image, image variation |
| [VQ Diffusion](vq_diffusion) | text2image |
| [Wuerstchen](wuerstchen) | text2image |
## DiffusionPipeline
[[autodoc]] DiffusionPipeline
@@ -92,7 +34,3 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
## FlaxDiffusionPipeline
[[autodoc]] pipelines.pipeline_flax_utils.FlaxDiffusionPipeline
## PushToHubMixin
[[autodoc]] utils.PushToHubMixin

View File

@@ -10,7 +10,7 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Paint By Example
# PaintByExample
[Paint by Example: Exemplar-based Image Editing with Diffusion Models](https://huggingface.co/papers/2211.13227) is by Binxin Yang, Shuyang Gu, Bo Zhang, Ting Zhang, Xuejin Chen, Xiaoyan Sun, Dong Chen, Fang Wen.
@@ -26,7 +26,7 @@ PaintByExample is supported by the official [Fantasy-Studio/Paint-by-Example](ht
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -44,7 +44,7 @@ But with circular padding, the right and the left parts are matching (`circular_
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -41,7 +41,7 @@ in parallel on multiple GPUs. But [`StableDiffusionParadigmsPipeline`] is design
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -22,7 +22,7 @@ You can find additional information about InstructPix2Pix on the [project page](
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
@@ -34,7 +34,5 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
- load_lora_weights
- save_lora_weights
## StableDiffusionXLInstructPix2PixPipeline
[[autodoc]] StableDiffusionXLInstructPix2PixPipeline
- __call__
- all
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput

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@@ -1,36 +0,0 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# PixArt
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/pixart/header_collage.png)
[PixArt-α: Fast Training of Diffusion Transformer for Photorealistic Text-to-Image Synthesis](https://huggingface.co/papers/2310.00426) is Junsong Chen, Jincheng Yu, Chongjian Ge, Lewei Yao, Enze Xie, Yue Wu, Zhongdao Wang, James Kwok, Ping Luo, Huchuan Lu, and Zhenguo Li.
The abstract from the paper is:
*The most advanced text-to-image (T2I) models require significant training costs (e.g., millions of GPU hours), seriously hindering the fundamental innovation for the AIGC community while increasing CO2 emissions. This paper introduces PIXART-α, a Transformer-based T2I diffusion model whose image generation quality is competitive with state-of-the-art image generators (e.g., Imagen, SDXL, and even Midjourney), reaching near-commercial application standards. Additionally, it supports high-resolution image synthesis up to 1024px resolution with low training cost, as shown in Figure 1 and 2. To achieve this goal, three core designs are proposed: (1) Training strategy decomposition: We devise three distinct training steps that separately optimize pixel dependency, text-image alignment, and image aesthetic quality; (2) Efficient T2I Transformer: We incorporate cross-attention modules into Diffusion Transformer (DiT) to inject text conditions and streamline the computation-intensive class-condition branch; (3) High-informative data: We emphasize the significance of concept density in text-image pairs and leverage a large Vision-Language model to auto-label dense pseudo-captions to assist text-image alignment learning. As a result, PIXART-α's training speed markedly surpasses existing large-scale T2I models, e.g., PIXART-α only takes 10.8% of Stable Diffusion v1.5's training time (675 vs. 6,250 A100 GPU days), saving nearly $300,000 ($26,000 vs. $320,000) and reducing 90% CO2 emissions. Moreover, compared with a larger SOTA model, RAPHAEL, our training cost is merely 1%. Extensive experiments demonstrate that PIXART-α excels in image quality, artistry, and semantic control. We hope PIXART-α will provide new insights to the AIGC community and startups to accelerate building their own high-quality yet low-cost generative models from scratch.*
You can find the original codebase at [PixArt-alpha/PixArt-alpha](https://github.com/PixArt-alpha/PixArt-alpha) and all the available checkpoints at [PixArt-alpha](https://huggingface.co/PixArt-alpha).
Some notes about this pipeline:
* It uses a Transformer backbone (instead of a UNet) for denoising. As such it has a similar architecture as [DiT](./dit.md).
* It was trained using text conditions computed from T5. This aspect makes the pipeline better at following complex text prompts with intricate details.
* It is good at producing high-resolution images at different aspect ratios. To get the best results, the authors recommend some size brackets which can be found [here](https://github.com/PixArt-alpha/PixArt-alpha/blob/08fbbd281ec96866109bdd2cdb75f2f58fb17610/diffusion/data/datasets/utils.py).
* It rivals the quality of state-of-the-art text-to-image generation systems (as of this writing) such as Stable Diffusion XL, Imagen, and DALL-E 2, while being more efficient than them.
## PixArtAlphaPipeline
[[autodoc]] PixArtAlphaPipeline
- all
- __call__

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@@ -22,7 +22,7 @@ The original codebase can be found at [luping-liu/PNDM](https://github.com/lupin
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

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@@ -23,7 +23,7 @@ The original codebase can be found at [andreas128/RePaint](https://github.com/an
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

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@@ -22,7 +22,7 @@ The original codebase can be found at [yang-song/score_sde_pytorch](https://gith
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -22,7 +22,7 @@ You can find additional information about Self-Attention Guidance on the [projec
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

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@@ -21,7 +21,7 @@ The abstract from the paper is:
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
@@ -31,5 +31,5 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
- __call__
## StableDiffusionSafePipelineOutput
[[autodoc]] pipelines.semantic_stable_diffusion.pipeline_output.SemanticStableDiffusionPipelineOutput
- all
[[autodoc]] pipelines.semantic_stable_diffusion.SemanticStableDiffusionPipelineOutput
- all

View File

@@ -9,7 +9,7 @@ specific language governing permissions and limitations under the License.
# Shap-E
The Shap-E model was proposed in [Shap-E: Generating Conditional 3D Implicit Functions](https://huggingface.co/papers/2305.02463) by Alex Nichol and Heewon Jun from [OpenAI](https://github.com/openai).
The Shap-E model was proposed in [Shap-E: Generating Conditional 3D Implicit Functions](https://huggingface.co/papers/2305.02463) by Alex Nichol and Heewon Jun from [OpenAI](https://github.com/openai).
The abstract from the paper is:
@@ -19,10 +19,163 @@ The original codebase can be found at [openai/shap-e](https://github.com/openai/
<Tip>
See the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Usage Examples
In the following, we will walk you through some examples of how to use Shap-E pipelines to create 3D objects in gif format.
### Text-to-3D image generation
We can use [`ShapEPipeline`] to create 3D object based on a text prompt. In this example, we will make a birthday cupcake for :firecracker: diffusers library's 1 year birthday. The workflow to use the Shap-E text-to-image pipeline is same as how you would use other text-to-image pipelines in diffusers.
```python
import torch
from diffusers import DiffusionPipeline
device = torch.device("cuda" if torch.cuda.is_available() else "cpu")
repo = "openai/shap-e"
pipe = DiffusionPipeline.from_pretrained(repo, torch_dtype=torch.float16)
pipe = pipe.to(device)
guidance_scale = 15.0
prompt = ["A firecracker", "A birthday cupcake"]
images = pipe(
prompt,
guidance_scale=guidance_scale,
num_inference_steps=64,
frame_size=256,
).images
```
The output of [`ShapEPipeline`] is a list of lists of images frames. Each list of frames can be used to create a 3D object. Let's use the `export_to_gif` utility function in diffusers to make a 3D cupcake!
```python
from diffusers.utils import export_to_gif
export_to_gif(images[0], "firecracker_3d.gif")
export_to_gif(images[1], "cake_3d.gif")
```
![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/shap_e/firecracker_out.gif)
![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/shap_e/cake_out.gif)
### Image-to-Image generation
You can use [`ShapEImg2ImgPipeline`] along with other text-to-image pipelines in diffusers and turn your 2D generation into 3D.
In this example, We will first genrate a cheeseburger with a simple prompt "A cheeseburger, white background"
```python
from diffusers import DiffusionPipeline
import torch
pipe_prior = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1-prior", torch_dtype=torch.float16)
pipe_prior.to("cuda")
t2i_pipe = DiffusionPipeline.from_pretrained("kandinsky-community/kandinsky-2-1", torch_dtype=torch.float16)
t2i_pipe.to("cuda")
prompt = "A cheeseburger, white background"
image_embeds, negative_image_embeds = pipe_prior(prompt, guidance_scale=1.0).to_tuple()
image = t2i_pipe(
prompt,
image_embeds=image_embeds,
negative_image_embeds=negative_image_embeds,
).images[0]
image.save("burger.png")
```
![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/shap_e/burger_in.png)
we will then use the Shap-E image-to-image pipeline to turn it into a 3D cheeseburger :)
```python
from PIL import Image
from diffusers.utils import export_to_gif
repo = "openai/shap-e-img2img"
pipe = DiffusionPipeline.from_pretrained(repo, torch_dtype=torch.float16)
pipe = pipe.to("cuda")
guidance_scale = 3.0
image = Image.open("burger.png").resize((256, 256))
images = pipe(
image,
guidance_scale=guidance_scale,
num_inference_steps=64,
frame_size=256,
).images
gif_path = export_to_gif(images[0], "burger_3d.gif")
```
![img](https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/shap_e/burger_out.gif)
### Generate mesh
For both [`ShapEPipeline`] and [`ShapEImg2ImgPipeline`], you can generate mesh output by passing `output_type` as `mesh` to the pipeline, and then use the [`ShapEPipeline.export_to_ply`] utility function to save the output as a `ply` file. We also provide a [`ShapEPipeline.export_to_obj`] function that you can use to save mesh outputs as `obj` files.
```python
import torch
from diffusers import DiffusionPipeline
from diffusers.utils import export_to_ply
device = torch.device("cuda" if torch.cuda.is_available() else "cpu")
repo = "openai/shap-e"
pipe = DiffusionPipeline.from_pretrained(repo, torch_dtype=torch.float16, variant="fp16")
pipe = pipe.to(device)
guidance_scale = 15.0
prompt = "A birthday cupcake"
images = pipe(prompt, guidance_scale=guidance_scale, num_inference_steps=64, frame_size=256, output_type="mesh").images
ply_path = export_to_ply(images[0], "3d_cake.ply")
print(f"saved to folder: {ply_path}")
```
Huggingface Datasets supports mesh visualization for mesh files in `glb` format. Below we will show you how to convert your mesh file into `glb` format so that you can use the Dataset viewer to render 3D objects.
We need to install `trimesh` library.
```
pip install trimesh
```
To convert the mesh file into `glb` format,
```python
import trimesh
mesh = trimesh.load("3d_cake.ply")
mesh.export("3d_cake.glb", file_type="glb")
```
By default, the mesh output of Shap-E is from the bottom viewpoint; you can change the default viewpoint by applying a rotation transformation
```python
import trimesh
import numpy as np
mesh = trimesh.load("3d_cake.ply")
rot = trimesh.transformations.rotation_matrix(-np.pi / 2, [1, 0, 0])
mesh = mesh.apply_transform(rot)
mesh.export("3d_cake.glb", file_type="glb")
```
Now you can upload your mesh file to your dataset and visualize it! Here is the link to the 3D cake we just generated
https://huggingface.co/datasets/hf-internal-testing/diffusers-images/blob/main/shap_e/3d_cake.glb
## ShapEPipeline
[[autodoc]] ShapEPipeline
- all

View File

@@ -24,7 +24,7 @@ As depicted above the model takes as input a MIDI file and tokenizes it into a s
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -28,12 +28,11 @@ This model was contributed by the community contributor [HimariO](https://github
| Pipeline | Tasks | Demo
|---|---|:---:|
| [StableDiffusionAdapterPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_adapter.py) | *Text-to-Image Generation with T2I-Adapter Conditioning* | -
| [StableDiffusionXLAdapterPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_xl_adapter.py) | *Text-to-Image Generation with T2I-Adapter Conditioning on StableDiffusion-XL* | -
| [StableDiffusionAdapterPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_adapter.py) | *Text-to-Image Generation with T2I-Adapter Conditioning* | -
## Usage example with the base model of StableDiffusion-1.4/1.5
## Usage example
In the following we give a simple example of how to use a *T2IAdapter* checkpoint with Diffusers for inference based on StableDiffusion-1.4/1.5.
In the following we give a simple example of how to use a *T2IAdapter* checkpoint with Diffusers for inference.
All adapters use the same pipeline.
1. Images are first converted into the appropriate *control image* format.
@@ -70,7 +69,7 @@ Next, create the adapter pipeline
import torch
from diffusers import StableDiffusionAdapterPipeline, T2IAdapter
adapter = T2IAdapter.from_pretrained("TencentARC/t2iadapter_color_sd14v1", torch_dtype=torch.float16)
adapter = T2IAdapter.from_pretrained("TencentARC/t2iadapter_color_sd14v1")
pipe = StableDiffusionAdapterPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
adapter=adapter,
@@ -94,62 +93,6 @@ out_image = pipe(
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_output.png)
## Usage example with the base model of StableDiffusion-XL
In the following we give a simple example of how to use a *T2IAdapter* checkpoint with Diffusers for inference based on StableDiffusion-XL.
All adapters use the same pipeline.
1. Images are first downloaded into the appropriate *control image* format.
2. The *control image* and *prompt* are passed to the [`StableDiffusionXLAdapterPipeline`].
Let's have a look at a simple example using the [Sketch Adapter](https://huggingface.co/Adapter/t2iadapter/tree/main/sketch_sdxl_1.0).
```python
from diffusers.utils import load_image
sketch_image = load_image("https://huggingface.co/Adapter/t2iadapter/resolve/main/sketch.png").convert("L")
```
![img](https://huggingface.co/Adapter/t2iadapter/resolve/main/sketch.png)
Then, create the adapter pipeline
```py
import torch
from diffusers import (
T2IAdapter,
StableDiffusionXLAdapterPipeline,
DDPMScheduler
)
from diffusers.models.unet_2d_condition import UNet2DConditionModel
model_id = "stabilityai/stable-diffusion-xl-base-1.0"
adapter = T2IAdapter.from_pretrained("Adapter/t2iadapter", subfolder="sketch_sdxl_1.0",torch_dtype=torch.float16, adapter_type="full_adapter_xl")
scheduler = DDPMScheduler.from_pretrained(model_id, subfolder="scheduler")
pipe = StableDiffusionXLAdapterPipeline.from_pretrained(
model_id, adapter=adapter, safety_checker=None, torch_dtype=torch.float16, variant="fp16", scheduler=scheduler
)
pipe.to("cuda")
```
Finally, pass the prompt and control image to the pipeline
```py
# fix the random seed, so you will get the same result as the example
generator = torch.Generator().manual_seed(42)
sketch_image_out = pipe(
prompt="a photo of a dog in real world, high quality",
negative_prompt="extra digit, fewer digits, cropped, worst quality, low quality",
image=sketch_image,
generator=generator,
guidance_scale=7.5
).images[0]
```
![img](https://huggingface.co/Adapter/t2iadapter/resolve/main/sketch_output.png)
## Available checkpoints
@@ -170,9 +113,6 @@ Non-diffusers checkpoints can be found under [TencentARC/T2I-Adapter](https://hu
|[TencentARC/t2iadapter_depth_sd15v2](https://huggingface.co/TencentARC/t2iadapter_depth_sd15v2)||
|[TencentARC/t2iadapter_sketch_sd15v2](https://huggingface.co/TencentARC/t2iadapter_sketch_sd15v2)||
|[TencentARC/t2iadapter_zoedepth_sd15v1](https://huggingface.co/TencentARC/t2iadapter_zoedepth_sd15v1)||
|[Adapter/t2iadapter, subfolder='sketch_sdxl_1.0'](https://huggingface.co/Adapter/t2iadapter/tree/main/sketch_sdxl_1.0)||
|[Adapter/t2iadapter, subfolder='canny_sdxl_1.0'](https://huggingface.co/Adapter/t2iadapter/tree/main/canny_sdxl_1.0)||
|[Adapter/t2iadapter, subfolder='openpose_sdxl_1.0'](https://huggingface.co/Adapter/t2iadapter/tree/main/openpose_sdxl_1.0)||
## Combining multiple adapters
@@ -245,14 +185,3 @@ However, T2I-Adapter performs slightly worse than ControlNet.
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
## StableDiffusionXLAdapterPipeline
[[autodoc]] StableDiffusionXLAdapterPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_vae_slicing
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention

View File

@@ -1,59 +0,0 @@
<!--Copyright 2023 The GLIGEN Authors and The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# GLIGEN (Grounded Language-to-Image Generation)
The GLIGEN model was created by researchers and engineers from [University of Wisconsin-Madison, Columbia University, and Microsoft](https://github.com/gligen/GLIGEN). The [`StableDiffusionGLIGENPipeline`] and [`StableDiffusionGLIGENTextImagePipeline`] can generate photorealistic images conditioned on grounding inputs. Along with text and bounding boxes with [`StableDiffusionGLIGENPipeline`], if input images are given, [`StableDiffusionGLIGENTextImagePipeline`] can insert objects described by text at the region defined by bounding boxes. Otherwise, it'll generate an image described by the caption/prompt and insert objects described by text at the region defined by bounding boxes. It's trained on COCO2014D and COCO2014CD datasets, and the model uses a frozen CLIP ViT-L/14 text encoder to condition itself on grounding inputs.
The abstract from the [paper](https://huggingface.co/papers/2301.07093) is:
*Large-scale text-to-image diffusion models have made amazing advances. However, the status quo is to use text input alone, which can impede controllability. In this work, we propose GLIGEN, Grounded-Language-to-Image Generation, a novel approach that builds upon and extends the functionality of existing pre-trained text-to-image diffusion models by enabling them to also be conditioned on grounding inputs. To preserve the vast concept knowledge of the pre-trained model, we freeze all of its weights and inject the grounding information into new trainable layers via a gated mechanism. Our model achieves open-world grounded text2img generation with caption and bounding box condition inputs, and the grounding ability generalizes well to novel spatial configurations and concepts. GLIGENs zeroshot performance on COCO and LVIS outperforms existing supervised layout-to-image baselines by a large margin.*
<Tip>
Make sure to check out the Stable Diffusion [Tips](https://huggingface.co/docs/diffusers/en/api/pipelines/stable_diffusion/overview#tips) section to learn how to explore the tradeoff between scheduler speed and quality and how to reuse pipeline components efficiently!
If you want to use one of the official checkpoints for a task, explore the [gligen](https://huggingface.co/gligen) Hub organizations!
</Tip>
[`StableDiffusionGLIGENPipeline`] was contributed by [Nikhil Gajendrakumar](https://github.com/nikhil-masterful) and [`StableDiffusionGLIGENTextImagePipeline`] was contributed by [Nguyễn Công Tú Anh](https://github.com/tuanh123789).
## StableDiffusionGLIGENPipeline
[[autodoc]] StableDiffusionGLIGENPipeline
- all
- __call__
- enable_vae_slicing
- disable_vae_slicing
- enable_vae_tiling
- disable_vae_tiling
- enable_model_cpu_offload
- prepare_latents
- enable_fuser
## StableDiffusionGLIGENTextImagePipeline
[[autodoc]] StableDiffusionGLIGENTextImagePipeline
- all
- __call__
- enable_vae_slicing
- disable_vae_slicing
- enable_vae_tiling
- disable_vae_tiling
- enable_model_cpu_offload
- prepare_latents
- enable_fuser
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput

View File

@@ -12,7 +12,7 @@ specific language governing permissions and limitations under the License.
# Text-to-(RGB, depth)
LDM3D was proposed in [LDM3D: Latent Diffusion Model for 3D](https://huggingface.co/papers/2305.10853) by Gabriela Ben Melech Stan, Diana Wofk, Scottie Fox, Alex Redden, Will Saxton, Jean Yu, Estelle Aflalo, Shao-Yen Tseng, Fabio Nonato, Matthias Muller, and Vasudev Lal. LDM3D generates an image and a depth map from a given text prompt unlike the existing text-to-image diffusion models such as [Stable Diffusion](./overview) which only generates an image. With almost the same number of parameters, LDM3D achieves to create a latent space that can compress both the RGB images and the depth maps.
LDM3D was proposed in [LDM3D: Latent Diffusion Model for 3D](https://huggingface.co/papers/2305.10853) by Gabriela Ben Melech Stan, Diana Wofk, Scottie Fox, Alex Redden, Will Saxton, Jean Yu, Estelle Aflalo, Shao-Yen Tseng, Fabio Nonato, Matthias Muller, and Vasudev Lal. LDM3D generates an image and a depth map from a given text prompt unlike the existing text-to-image diffusion models such as [Stable Diffusion](./stable_diffusion/overview) which only generates an image. With almost the same number of parameters, LDM3D achieves to create a latent space that can compress both the RGB images and the depth maps.
The abstract from the paper is:
@@ -30,8 +30,8 @@ Make sure to check out the Stable Diffusion [Tips](overview#tips) section to lea
- all
- __call__
## LDM3DPipelineOutput
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.pipeline_stable_diffusion_ldm3d.LDM3DPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
- all
- __call__

View File

@@ -36,8 +36,10 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<th class="px-4 py-2 font-medium text-gray-900 text-left">
Space
</th>
</tr>
</thead>
<tbody class="divide-y divide-gray-200">
<tr>
<td class="px-4 py-2 text-gray-700">
@@ -47,6 +49,7 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/stabilityai/stable-diffusion"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./img2img">StableDiffusionImg2Img</a>
@@ -55,6 +58,7 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/huggingface/diffuse-the-rest"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./inpaint">StableDiffusionInpaint</a>
@@ -63,6 +67,7 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/runwayml/stable-diffusion-inpainting"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./depth2img">StableDiffusionDepth2Img</a>
@@ -71,6 +76,7 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/radames/stable-diffusion-depth2img"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./image_variation">StableDiffusionImageVariation</a>
@@ -79,6 +85,7 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/lambdalabs/stable-diffusion-image-variations"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./stable_diffusion_safe">StableDiffusionPipelineSafe</a>
@@ -87,6 +94,7 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/AIML-TUDA/unsafe-vs-safe-stable-diffusion"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./stable_diffusion_2">StableDiffusion2</a>
@@ -95,6 +103,7 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/stabilityai/stable-diffusion"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./stable_diffusion_xl">StableDiffusionXL</a>
@@ -103,6 +112,7 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/RamAnanth1/stable-diffusion-xl"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./latent_upscale">StableDiffusionLatentUpscale</a>
@@ -111,12 +121,14 @@ The table below summarizes the available Stable Diffusion pipelines, their suppo
<td class="px-4 py-2"><a href="https://huggingface.co/spaces/huggingface-projects/stable-diffusion-latent-upscaler"><img src="https://img.shields.io/badge/%F0%9F%A4%97%20Hugging%20Face-Spaces-blue"/></a>
</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./upscale">StableDiffusionUpscale</a>
</td>
<td class="px-4 py-2 text-gray-700">super-resolution</td>
</tr>
<tr>
<td class="px-4 py-2 text-gray-700">
<a href="./ldm3d_diffusion">StableDiffusionLDM3D</a>

View File

@@ -10,32 +10,382 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Stable Diffusion XL
# Stable diffusion XL
Stable Diffusion XL (SDXL) was proposed in [SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis](https://huggingface.co/papers/2307.01952) by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, and Robin Rombach.
Stable Diffusion XL was proposed in [SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis](https://arxiv.org/abs/2307.01952) by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, Robin Rombach
The abstract from the paper is:
The abstract of the paper is the following:
*We present SDXL, a latent diffusion model for text-to-image synthesis. Compared to previous versions of Stable Diffusion, SDXL leverages a three times larger UNet backbone: The increase of model parameters is mainly due to more attention blocks and a larger cross-attention context as SDXL uses a second text encoder. We design multiple novel conditioning schemes and train SDXL on multiple aspect ratios. We also introduce a refinement model which is used to improve the visual fidelity of samples generated by SDXL using a post-hoc image-to-image technique. We demonstrate that SDXL shows drastically improved performance compared the previous versions of Stable Diffusion and achieves results competitive with those of black-box state-of-the-art image generators.*
## Tips
- Using SDXL with a DPM++ scheduler for less than 50 steps is known to produce [visual artifacts](https://github.com/huggingface/diffusers/issues/5433) because the solver becomes numerically unstable. To fix this issue, take a look at this [PR](https://github.com/huggingface/diffusers/pull/5541) which recommends for ODE/SDE solvers:
- set `use_karras_sigmas=True` or `lu_lambdas=True` to improve image quality
- set `euler_at_final=True` if you're using a solver with uniform step sizes (DPM++2M or DPM++2M SDE)
- Most SDXL checkpoints work best with an image size of 1024x1024. Image sizes of 768x768 and 512x512 are also supported, but the results aren't as good. Anything below 512x512 is not recommended and likely won't for for default checkpoints like [stabilityai/stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0).
- SDXL can pass a different prompt for each of the text encoders it was trained on. We can even pass different parts of the same prompt to the text encoders.
- SDXL output images can be improved by making use of a refiner model in an image-to-image setting.
- SDXL offers `negative_original_size`, `negative_crops_coords_top_left`, and `negative_target_size` to negatively condition the model on image resolution and cropping parameters.
- Stable Diffusion XL works especially well with images between 768 and 1024.
- Stable Diffusion XL can pass a different prompt for each of the text encoders it was trained on as shown below. We can even pass different parts of the same prompt to the text encoders.
- Stable Diffusion XL output image can be improved by making use of a refiner as shown below.
### Available checkpoints:
- *Text-to-Image (1024x1024 resolution)*: [stabilityai/stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) with [`StableDiffusionXLPipeline`]
- *Image-to-Image / Refiner (1024x1024 resolution)*: [stabilityai/stable-diffusion-xl-refiner-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0) with [`StableDiffusionXLImg2ImgPipeline`]
## Usage Example
Before using SDXL make sure to have `transformers`, `accelerate`, `safetensors` and `invisible_watermark` installed.
You can install the libraries as follows:
```
pip install transformers
pip install accelerate
pip install safetensors
```
### Watermarker
We recommend to add an invisible watermark to images generating by Stable Diffusion XL, this can help with identifying if an image is machine-synthesised for downstream applications. To do so, please install
the [invisible-watermark library](https://pypi.org/project/invisible-watermark/) via:
```
pip install invisible-watermark>=0.2.0
```
If the `invisible-watermark` library is installed the watermarker will be used **by default**.
If you have other provisions for generating or deploying images safely, you can disable the watermarker as follows:
```py
pipe = StableDiffusionXLPipeline.from_pretrained(..., add_watermarker=False)
```
### Text-to-Image
You can use SDXL as follows for *text-to-image*:
```py
from diffusers import StableDiffusionXLPipeline
import torch
pipe = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
)
pipe.to("cuda")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
image = pipe(prompt=prompt).images[0]
```
### Image-to-image
You can use SDXL as follows for *image-to-image*:
```py
import torch
from diffusers import StableDiffusionXLImg2ImgPipeline
from diffusers.utils import load_image
pipe = StableDiffusionXLImg2ImgPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-refiner-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
)
pipe = pipe.to("cuda")
url = "https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/aa_xl/000000009.png"
init_image = load_image(url).convert("RGB")
prompt = "a photo of an astronaut riding a horse on mars"
image = pipe(prompt, image=init_image).images[0]
```
### Inpainting
You can use SDXL as follows for *inpainting*
```py
import torch
from diffusers import StableDiffusionXLInpaintPipeline
from diffusers.utils import load_image
pipe = StableDiffusionXLInpaintPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
)
pipe.to("cuda")
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
init_image = load_image(img_url).convert("RGB")
mask_image = load_image(mask_url).convert("RGB")
prompt = "A majestic tiger sitting on a bench"
image = pipe(prompt=prompt, image=init_image, mask_image=mask_image, num_inference_steps=50, strength=0.80).images[0]
```
### Refining the image output
In addition to the [base model checkpoint](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0),
StableDiffusion-XL also includes a [refiner checkpoint](huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0)
that is specialized in denoising low-noise stage images to generate images of improved high-frequency quality.
This refiner checkpoint can be used as a "second-step" pipeline after having run the base checkpoint to improve
image quality.
When using the refiner, one can easily
- 1.) employ the base model and refiner as an *Ensemble of Expert Denoisers* as first proposed in [eDiff-I](https://research.nvidia.com/labs/dir/eDiff-I/) or
- 2.) simply run the refiner in [SDEdit](https://arxiv.org/abs/2108.01073) fashion after the base model.
**Note**: The idea of using SD-XL base & refiner as an ensemble of experts was first brought forward by
a couple community contributors which also helped shape the following `diffusers` implementation, namely:
- [SytanSD](https://github.com/SytanSD)
- [bghira](https://github.com/bghira)
- [Birch-san](https://github.com/Birch-san)
- [AmericanPresidentJimmyCarter](https://github.com/AmericanPresidentJimmyCarter)
#### 1.) Ensemble of Expert Denoisers
When using the base and refiner model as an ensemble of expert of denoisers, the base model should serve as the
expert for the high-noise diffusion stage and the refiner serves as the expert for the low-noise diffusion stage.
The advantage of 1.) over 2.) is that it requires less overall denoising steps and therefore should be significantly
faster. The drawback is that one cannot really inspect the output of the base model; it will still be heavily denoised.
To use the base model and refiner as an ensemble of expert denoisers, make sure to define the span
of timesteps which should be run through the high-noise denoising stage (*i.e.* the base model) and the low-noise
denoising stage (*i.e.* the refiner model) respectively. We can set the intervals using the [`denoising_end`](https://huggingface.co/docs/diffusers/main/en/api/pipelines/stable_diffusion/stable_diffusion_xl#diffusers.StableDiffusionXLPipeline.__call__.denoising_end) of the base model
and [`denoising_start`](https://huggingface.co/docs/diffusers/main/en/api/pipelines/stable_diffusion/stable_diffusion_xl#diffusers.StableDiffusionXLImg2ImgPipeline.__call__.denoising_start) of the refiner model.
For both `denoising_end` and `denoising_start` a float value between 0 and 1 should be passed.
When passed, the end and start of denoising will be defined by proportions of discrete timesteps as
defined by the model schedule.
Note that this will override `strength` if it is also declared, since the number of denoising steps
is determined by the discrete timesteps the model was trained on and the declared fractional cutoff.
Let's look at an example.
First, we import the two pipelines. Since the text encoders and variational autoencoder are the same
you don't have to load those again for the refiner.
```py
from diffusers import DiffusionPipeline
import torch
base = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
)
base.to("cuda")
refiner = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-refiner-1.0",
text_encoder_2=base.text_encoder_2,
vae=base.vae,
torch_dtype=torch.float16,
use_safetensors=True,
variant="fp16",
)
refiner.to("cuda")
```
Now we define the number of inference steps and the point at which the model shall be run through the
high-noise denoising stage (*i.e.* the base model).
```py
n_steps = 40
high_noise_frac = 0.8
```
Stable Diffusion XL base is trained on timesteps 0-999 and Stable Diffusion XL refiner is finetuned
from the base model on low noise timesteps 0-199 inclusive, so we use the base model for the first
800 timesteps (high noise) and the refiner for the last 200 timesteps (low noise). Hence, `high_noise_frac`
is set to 0.8, so that all steps 200-999 (the first 80% of denoising timesteps) are performed by the
base model and steps 0-199 (the last 20% of denoising timesteps) are performed by the refiner model.
Remember, the denoising process starts at **high value** (high noise) timesteps and ends at
**low value** (low noise) timesteps.
Let's run the two pipelines now. Make sure to set `denoising_end` and
`denoising_start` to the same values and keep `num_inference_steps` constant. Also remember that
the output of the base model should be in latent space:
```py
prompt = "A majestic lion jumping from a big stone at night"
image = base(
prompt=prompt,
num_inference_steps=n_steps,
denoising_end=high_noise_frac,
output_type="latent",
).images
image = refiner(
prompt=prompt,
num_inference_steps=n_steps,
denoising_start=high_noise_frac,
image=image,
).images[0]
```
Let's have a look at the images
| Original Image | Ensemble of Denoisers Experts |
|---|---|
| ![lion_base_timesteps](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lion_base.png) | ![lion_refined_timesteps](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lion_refined.png)
If we would have just run the base model on the same 40 steps, the image would have been arguably less detailed (e.g. the lion eyes and nose):
<Tip>
To learn how to use SDXL for various tasks, how to optimize performance, and other usage examples, take a look at the [Stable Diffusion XL](../../../using-diffusers/sdxl) guide.
Check out the [Stability AI](https://huggingface.co/stabilityai) Hub organization for the official base and refiner model checkpoints!
The ensemble-of-experts method works well on all available schedulers!
</Tip>
#### 2.) Refining the image output from fully denoised base image
In standard [`StableDiffusionImg2ImgPipeline`]-fashion, the fully-denoised image generated of the base model
can be further improved using the [refiner checkpoint](huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0).
For this, you simply run the refiner as a normal image-to-image pipeline after the "base" text-to-image
pipeline. You can leave the outputs of the base model in latent space.
```py
from diffusers import DiffusionPipeline
import torch
pipe = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
)
pipe.to("cuda")
refiner = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-refiner-1.0",
text_encoder_2=pipe.text_encoder_2,
vae=pipe.vae,
torch_dtype=torch.float16,
use_safetensors=True,
variant="fp16",
)
refiner.to("cuda")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
image = pipe(prompt=prompt, output_type="latent" if use_refiner else "pil").images[0]
image = refiner(prompt=prompt, image=image[None, :]).images[0]
```
| Original Image | Refined Image |
|---|---|
| ![](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/sd_xl/init_image.png) | ![](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/sd_xl/refined_image.png) |
<Tip>
The refiner can also very well be used in an in-painting setting. To do so just make
sure you use the [`StableDiffusionXLInpaintPipeline`] classes as shown below
</Tip>
To use the refiner for inpainting in the Ensemble of Expert Denoisers setting you can do the following:
```py
from diffusers import StableDiffusionXLInpaintPipeline
from diffusers.utils import load_image
pipe = StableDiffusionXLInpaintPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
)
pipe.to("cuda")
refiner = StableDiffusionXLInpaintPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-refiner-1.0",
text_encoder_2=pipe.text_encoder_2,
vae=pipe.vae,
torch_dtype=torch.float16,
use_safetensors=True,
variant="fp16",
)
refiner.to("cuda")
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
init_image = load_image(img_url).convert("RGB")
mask_image = load_image(mask_url).convert("RGB")
prompt = "A majestic tiger sitting on a bench"
num_inference_steps = 75
high_noise_frac = 0.7
image = pipe(
prompt=prompt,
image=init_image,
mask_image=mask_image,
num_inference_steps=num_inference_steps,
denoising_start=high_noise_frac,
output_type="latent",
).images
image = refiner(
prompt=prompt,
image=image,
mask_image=mask_image,
num_inference_steps=num_inference_steps,
denoising_start=high_noise_frac,
).images[0]
```
To use the refiner for inpainting in the standard SDE-style setting, simply remove `denoising_end` and `denoising_start` and choose a smaller
number of inference steps for the refiner.
### Loading single file checkpoints / original file format
By making use of [`~diffusers.loaders.FromSingleFileMixin.from_single_file`] you can also load the
original file format into `diffusers`:
```py
from diffusers import StableDiffusionXLPipeline, StableDiffusionXLImg2ImgPipeline
import torch
pipe = StableDiffusionXLPipeline.from_single_file(
"./sd_xl_base_1.0.safetensors", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
)
pipe.to("cuda")
refiner = StableDiffusionXLImg2ImgPipeline.from_single_file(
"./sd_xl_refiner_1.0.safetensors", torch_dtype=torch.float16, use_safetensors=True, variant="fp16"
)
refiner.to("cuda")
```
### Memory optimization via model offloading
If you are seeing out-of-memory errors, we recommend making use of [`StableDiffusionXLPipeline.enable_model_cpu_offload`].
```diff
- pipe.to("cuda")
+ pipe.enable_model_cpu_offload()
```
and
```diff
- refiner.to("cuda")
+ refiner.enable_model_cpu_offload()
```
### Speed-up inference with `torch.compile`
You can speed up inference by making use of `torch.compile`. This should give you **ca.** 20% speed-up.
```diff
+ pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
+ refiner.unet = torch.compile(refiner.unet, mode="reduce-overhead", fullgraph=True)
```
### Running with `torch < 2.0`
**Note** that if you want to run Stable Diffusion XL with `torch` < 2.0, please make sure to enable xformers
attention:
```
pip install xformers
```
```diff
+pipe.enable_xformers_memory_efficient_attention()
+refiner.enable_xformers_memory_efficient_attention()
```
## StableDiffusionXLPipeline
[[autodoc]] StableDiffusionXLPipeline
@@ -53,3 +403,25 @@ Check out the [Stability AI](https://huggingface.co/stabilityai) Hub organizatio
[[autodoc]] StableDiffusionXLInpaintPipeline
- all
- __call__
### Passing different prompts to each text-encoder
Stable Diffusion XL was trained on two text encoders. The default behavior is to pass the same prompt to each. But it is possible to pass a different prompt for each text-encoder, as [some users](https://github.com/huggingface/diffusers/issues/4004#issuecomment-1627764201) noted that it can boost quality.
To do so, you can pass `prompt_2` and `negative_prompt_2` in addition to `prompt` and `negative_prompt`. By doing that, you will pass the original prompts and negative prompts (as in `prompt` and `negative_prompt`) to `text_encoder` (in official SDXL 0.9/1.0 that is [OpenAI CLIP-ViT/L-14](https://huggingface.co/openai/clip-vit-large-patch14)),
and `prompt_2` and `negative_prompt_2` to `text_encoder_2` (in official SDXL 0.9/1.0 that is [OpenCLIP-ViT/bigG-14](https://huggingface.co/laion/CLIP-ViT-bigG-14-laion2B-39B-b160k)).
```py
from diffusers import StableDiffusionXLPipeline
import torch
pipe = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
)
pipe.to("cuda")
# prompt will be passed to OAI CLIP-ViT/L-14
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
# prompt_2 will be passed to OpenCLIP-ViT/bigG-14
prompt_2 = "monet painting"
image = pipe(prompt=prompt, prompt_2=prompt_2).images[0]
```

View File

@@ -20,7 +20,7 @@ The abstract from the paper:
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -7,9 +7,9 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# unCLIP
# UnCLIP
[Hierarchical Text-Conditional Image Generation with CLIP Latents](https://huggingface.co/papers/2204.06125) is by Aditya Ramesh, Prafulla Dhariwal, Alex Nichol, Casey Chu, Mark Chen. The unCLIP model in 🤗 Diffusers comes from kakaobrain's [karlo]((https://github.com/kakaobrain/karlo)).
[Hierarchical Text-Conditional Image Generation with CLIP Latents](https://huggingface.co/papers/2204.06125) is by Aditya Ramesh, Prafulla Dhariwal, Alex Nichol, Casey Chu, Mark Chen. The UnCLIP model in 🤗 Diffusers comes from kakaobrain's [karlo]((https://github.com/kakaobrain/karlo)).
The abstract from the paper is following:
@@ -19,7 +19,7 @@ You can find lucidrains DALL-E 2 recreation at [lucidrains/DALLE2-pytorch](https
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
@@ -34,4 +34,4 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
- __call__
## ImagePipelineOutput
[[autodoc]] pipelines.ImagePipelineOutput
[[autodoc]] pipelines.ImagePipelineOutput

View File

@@ -20,12 +20,6 @@ The abstract from the [paper](https://arxiv.org/abs/2303.06555) is:
You can find the original codebase at [thu-ml/unidiffuser](https://github.com/thu-ml/unidiffuser) and additional checkpoints at [thu-ml](https://huggingface.co/thu-ml).
<Tip warning={true}>
There is currently an issue on PyTorch 1.X where the output images are all black or the pixel values become `NaNs`. This issue can be mitigated by switching to PyTorch 2.X.
</Tip>
This pipeline was contributed by [dg845](https://github.com/dg845). ❤️
## Usage Examples

View File

@@ -31,7 +31,7 @@ You can load the more memory intensive "all-in-one" [`VersatileDiffusionPipeline
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -22,7 +22,7 @@ The original codebase can be found at [microsoft/VQ-Diffusion](https://github.co
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](/using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](/using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>

View File

@@ -1,149 +0,0 @@
# Würstchen
<img src="https://github.com/dome272/Wuerstchen/assets/61938694/0617c863-165a-43ee-9303-2a17299a0cf9">
[Würstchen: Efficient Pretraining of Text-to-Image Models](https://huggingface.co/papers/2306.00637) is by Pablo Pernias, Dominic Rampas, Mats L. Richter and Christopher Pal and Marc Aubreville.
The abstract from the paper is:
*We introduce Würstchen, a novel technique for text-to-image synthesis that unites competitive performance with unprecedented cost-effectiveness and ease of training on constrained hardware. Building on recent advancements in machine learning, our approach, which utilizes latent diffusion strategies at strong latent image compression rates, significantly reduces the computational burden, typically associated with state-of-the-art models, while preserving, if not enhancing, the quality of generated images. Wuerstchen achieves notable speed improvements at inference time, thereby rendering real-time applications more viable. One of the key advantages of our method lies in its modest training requirements of only 9,200 GPU hours, slashing the usual costs significantly without compromising the end performance. In a comparison against the state-of-the-art, we found the approach to yield strong competitiveness. This paper opens the door to a new line of research that prioritizes both performance and computational accessibility, hence democratizing the use of sophisticated AI technologies. Through Wuerstchen, we demonstrate a compelling stride forward in the realm of text-to-image synthesis, offering an innovative path to explore in future research.*
## Würstchen Overview
Würstchen is a diffusion model, whose text-conditional model works in a highly compressed latent space of images. Why is this important? Compressing data can reduce computational costs for both training and inference by magnitudes. Training on 1024x1024 images is way more expensive than training on 32x32. Usually, other works make use of a relatively small compression, in the range of 4x - 8x spatial compression. Würstchen takes this to an extreme. Through its novel design, we achieve a 42x spatial compression. This was unseen before because common methods fail to faithfully reconstruct detailed images after 16x spatial compression. Würstchen employs a two-stage compression, what we call Stage A and Stage B. Stage A is a VQGAN, and Stage B is a Diffusion Autoencoder (more details can be found in the [paper](https://huggingface.co/papers/2306.00637) ). A third model, Stage C, is learned in that highly compressed latent space. This training requires fractions of the compute used for current top-performing models, while also allowing cheaper and faster inference.
## Würstchen v2 comes to Diffusers
After the initial paper release, we have improved numerous things in the architecture, training and sampling, making Würstchen competitive to current state-of-the-art models in many ways. We are excited to release this new version together with Diffusers. Here is a list of the improvements.
- Higher resolution (1024x1024 up to 2048x2048)
- Faster inference
- Multi Aspect Resolution Sampling
- Better quality
We are releasing 3 checkpoints for the text-conditional image generation model (Stage C). Those are:
- v2-base
- v2-aesthetic
- **(default)** v2-interpolated (50% interpolation between v2-base and v2-aesthetic)
We recommend using v2-interpolated, as it has a nice touch of both photorealism and aesthetics. Use v2-base for finetunings as it does not have a style bias and use v2-aesthetic for very artistic generations.
A comparison can be seen here:
<img src="https://github.com/dome272/Wuerstchen/assets/61938694/2914830f-cbd3-461c-be64-d50734f4b49d" width=500>
## Text-to-Image Generation
For the sake of usability, Würstchen can be used with a single pipeline. This pipeline can be used as follows:
```python
import torch
from diffusers import AutoPipelineForText2Image
from diffusers.pipelines.wuerstchen import DEFAULT_STAGE_C_TIMESTEPS
pipe = AutoPipelineForText2Image.from_pretrained("warp-ai/wuerstchen", torch_dtype=torch.float16).to("cuda")
caption = "Anthropomorphic cat dressed as a fire fighter"
images = pipe(
caption,
width=1024,
height=1536,
prior_timesteps=DEFAULT_STAGE_C_TIMESTEPS,
prior_guidance_scale=4.0,
num_images_per_prompt=2,
).images
```
For explanation purposes, we can also initialize the two main pipelines of Würstchen individually. Würstchen consists of 3 stages: Stage C, Stage B, Stage A. They all have different jobs and work only together. When generating text-conditional images, Stage C will first generate the latents in a very compressed latent space. This is what happens in the `prior_pipeline`. Afterwards, the generated latents will be passed to Stage B, which decompresses the latents into a bigger latent space of a VQGAN. These latents can then be decoded by Stage A, which is a VQGAN, into the pixel-space. Stage B & Stage A are both encapsulated in the `decoder_pipeline`. For more details, take a look at the [paper](https://huggingface.co/papers/2306.00637).
```python
import torch
from diffusers import WuerstchenDecoderPipeline, WuerstchenPriorPipeline
from diffusers.pipelines.wuerstchen import DEFAULT_STAGE_C_TIMESTEPS
device = "cuda"
dtype = torch.float16
num_images_per_prompt = 2
prior_pipeline = WuerstchenPriorPipeline.from_pretrained(
"warp-ai/wuerstchen-prior", torch_dtype=dtype
).to(device)
decoder_pipeline = WuerstchenDecoderPipeline.from_pretrained(
"warp-ai/wuerstchen", torch_dtype=dtype
).to(device)
caption = "Anthropomorphic cat dressed as a fire fighter"
negative_prompt = ""
prior_output = prior_pipeline(
prompt=caption,
height=1024,
width=1536,
timesteps=DEFAULT_STAGE_C_TIMESTEPS,
negative_prompt=negative_prompt,
guidance_scale=4.0,
num_images_per_prompt=num_images_per_prompt,
)
decoder_output = decoder_pipeline(
image_embeddings=prior_output.image_embeddings,
prompt=caption,
negative_prompt=negative_prompt,
guidance_scale=0.0,
output_type="pil",
).images
```
## Speed-Up Inference
You can make use of `torch.compile` function and gain a speed-up of about 2-3x:
```python
prior_pipeline.prior = torch.compile(prior_pipeline.prior, mode="reduce-overhead", fullgraph=True)
decoder_pipeline.decoder = torch.compile(decoder_pipeline.decoder, mode="reduce-overhead", fullgraph=True)
```
## Limitations
- Due to the high compression employed by Würstchen, generations can lack a good amount
of detail. To our human eye, this is especially noticeable in faces, hands etc.
- **Images can only be generated in 128-pixel steps**, e.g. the next higher resolution
after 1024x1024 is 1152x1152
- The model lacks the ability to render correct text in images
- The model often does not achieve photorealism
- Difficult compositional prompts are hard for the model
The original codebase, as well as experimental ideas, can be found at [dome272/Wuerstchen](https://github.com/dome272/Wuerstchen).
## WuerstchenCombinedPipeline
[[autodoc]] WuerstchenCombinedPipeline
- all
- __call__
## WuerstchenPriorPipeline
[[autodoc]] WuerstchenPriorPipeline
- all
- __call__
## WuerstchenPriorPipelineOutput
[[autodoc]] pipelines.wuerstchen.pipeline_wuerstchen_prior.WuerstchenPriorPipelineOutput
## WuerstchenDecoderPipeline
[[autodoc]] WuerstchenDecoderPipeline
- all
- __call__
## Citation
```bibtex
@misc{pernias2023wuerstchen,
title={Wuerstchen: Efficient Pretraining of Text-to-Image Models},
author={Pablo Pernias and Dominic Rampas and Mats L. Richter and Christopher Pal and Marc Aubreville},
year={2023},
eprint={2306.00637},
archivePrefix={arXiv},
primaryClass={cs.CV}
}
```

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@@ -1,15 +1,11 @@
# CMStochasticIterativeScheduler
# Consistency Model Multistep Scheduler
[Consistency Models](https://huggingface.co/papers/2303.01469) by Yang Song, Prafulla Dhariwal, Mark Chen, and Ilya Sutskever introduced a multistep and onestep scheduler (Algorithm 1) that is capable of generating good samples in one or a small number of steps.
## Overview
The abstract from the paper is:
*Diffusion models have made significant breakthroughs in image, audio, and video generation, but they depend on an iterative generation process that causes slow sampling speed and caps their potential for real-time applications. To overcome this limitation, we propose consistency models, a new family of generative models that achieve high sample quality without adversarial training. They support fast one-step generation by design, while still allowing for few-step sampling to trade compute for sample quality. They also support zero-shot data editing, like image inpainting, colorization, and super-resolution, without requiring explicit training on these tasks. Consistency models can be trained either as a way to distill pre-trained diffusion models, or as standalone generative models. Through extensive experiments, we demonstrate that they outperform existing distillation techniques for diffusion models in one- and few-step generation. For example, we achieve the new state-of-the-art FID of 3.55 on CIFAR-10 and 6.20 on ImageNet 64x64 for one-step generation. When trained as standalone generative models, consistency models also outperform single-step, non-adversarial generative models on standard benchmarks like CIFAR-10, ImageNet 64x64 and LSUN 256x256.*
The original codebase can be found at [openai/consistency_models](https://github.com/openai/consistency_models).
Multistep and onestep scheduler (Algorithm 1) introduced alongside consistency models in the paper [Consistency Models](https://arxiv.org/abs/2303.01469) by Yang Song, Prafulla Dhariwal, Mark Chen, and Ilya Sutskever.
Based on the [original consistency models implementation](https://github.com/openai/consistency_models).
Should generate good samples from [`ConsistencyModelPipeline`] in one or a small number of steps.
## CMStochasticIterativeScheduler
[[autodoc]] CMStochasticIterativeScheduler
## CMStochasticIterativeSchedulerOutput
[[autodoc]] schedulers.scheduling_consistency_models.CMStochasticIterativeSchedulerOutput

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# ConsistencyDecoderScheduler
This scheduler is a part of the [`ConsistencyDecoderPipeline`] and was introduced in [DALL-E 3](https://openai.com/dall-e-3).
The original codebase can be found at [openai/consistency_models](https://github.com/openai/consistency_models).
## ConsistencyDecoderScheduler
[[autodoc]] schedulers.scheduling_consistency_decoder.ConsistencyDecoderScheduler

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specific language governing permissions and limitations under the License.
-->
# DDIMScheduler
# Denoising Diffusion Implicit Models (DDIM)
[Denoising Diffusion Implicit Models](https://huggingface.co/papers/2010.02502) (DDIM) by Jiaming Song, Chenlin Meng and Stefano Ermon.
## Overview
The abstract from the paper is:
[Denoising Diffusion Implicit Models](https://arxiv.org/abs/2010.02502) (DDIM) by Jiaming Song, Chenlin Meng and Stefano Ermon.
The abstract of the paper is the following:
*Denoising diffusion probabilistic models (DDPMs) have achieved high quality image generation without adversarial training,
yet they require simulating a Markov chain for many steps to produce a sample.
@@ -24,43 +26,50 @@ We construct a class of non-Markovian diffusion processes that lead to the same
We empirically demonstrate that DDIMs can produce high quality samples 10× to 50× faster in terms of wall-clock time compared to DDPMs, allow us to trade off
computation for sample quality, and can perform semantically meaningful image interpolation directly in the latent space.*
The original codebase of this paper can be found at [ermongroup/ddim](https://github.com/ermongroup/ddim), and you can contact the author on [tsong.me](https://tsong.me/).
The original codebase of this paper can be found here: [ermongroup/ddim](https://github.com/ermongroup/ddim).
For questions, feel free to contact the author on [tsong.me](https://tsong.me/).
## Tips
### Experimental: "Common Diffusion Noise Schedules and Sample Steps are Flawed":
The paper [Common Diffusion Noise Schedules and Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) claims that a mismatch between the training and inference settings leads to suboptimal inference generation results for Stable Diffusion. To fix this, the authors propose:
The paper **[Common Diffusion Noise Schedules and Sample Steps are Flawed](https://arxiv.org/abs/2305.08891)**
claims that a mismatch between the training and inference settings leads to suboptimal inference generation results for Stable Diffusion.
<Tip warning={true}>
The abstract reads as follows:
🧪 This is an experimental feature!
</Tip>
1. rescale the noise schedule to enforce zero terminal signal-to-noise ratio (SNR)
*We discover that common diffusion noise schedules do not enforce the last timestep to have zero signal-to-noise ratio (SNR),
and some implementations of diffusion samplers do not start from the last timestep.
Such designs are flawed and do not reflect the fact that the model is given pure Gaussian noise at inference, creating a discrepancy between training and inference.
We show that the flawed design causes real problems in existing implementations.
In Stable Diffusion, it severely limits the model to only generate images with medium brightness and
prevents it from generating very bright and dark samples. We propose a few simple fixes:
- (1) rescale the noise schedule to enforce zero terminal SNR;
- (2) train the model with v prediction;
- (3) change the sampler to always start from the last timestep;
- (4) rescale classifier-free guidance to prevent over-exposure.
These simple changes ensure the diffusion process is congruent between training and inference and
allow the model to generate samples more faithful to the original data distribution.*
You can apply all of these changes in `diffusers` when using [`DDIMScheduler`]:
- (1) rescale the noise schedule to enforce zero terminal SNR;
```py
pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config, rescale_betas_zero_snr=True)
```
2. train a model with `v_prediction` (add the following argument to the [train_text_to_image.py](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image.py) or [train_text_to_image_lora.py](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image_lora.py) scripts)
```bash
--prediction_type="v_prediction"
```
3. change the sampler to always start from the last timestep
- (2) train the model with v prediction;
Continue fine-tuning a checkpoint with [`train_text_to_image.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image.py) or [`train_text_to_image_lora.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image_lora.py)
and `--prediction_type="v_prediction"`.
- (3) change the sampler to always start from the last timestep;
```py
pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config, timestep_spacing="trailing")
```
4. rescale classifier-free guidance to prevent over-exposure
- (4) rescale classifier-free guidance to prevent over-exposure.
```py
image = pipeline(prompt, guidance_rescale=0.7).images[0]
pipe(..., guidance_rescale=0.7)
```
For example:
An example is to use [this checkpoint](https://huggingface.co/ptx0/pseudo-journey-v2)
which has been fine-tuned using the `"v_prediction"`.
The checkpoint can then be run in inference as follows:
```py
from diffusers import DiffusionPipeline, DDIMScheduler
@@ -77,6 +86,3 @@ image = pipeline(prompt, guidance_rescale=0.7).images[0]
## DDIMScheduler
[[autodoc]] DDIMScheduler
## DDIMSchedulerOutput
[[autodoc]] schedulers.scheduling_ddim.DDIMSchedulerOutput

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specific language governing permissions and limitations under the License.
-->
# DDIMInverseScheduler
# Inverse Denoising Diffusion Implicit Models (DDIMInverse)
`DDIMInverseScheduler` is the inverted scheduler from [Denoising Diffusion Implicit Models](https://huggingface.co/papers/2010.02502) (DDIM) by Jiaming Song, Chenlin Meng and Stefano Ermon.
The implementation is mostly based on the DDIM inversion definition from [Null-text Inversion for Editing Real Images using Guided Diffusion Models](https://huggingface.co/papers/2211.09794.pdf).
## Overview
This scheduler is the inverted scheduler of [Denoising Diffusion Implicit Models](https://arxiv.org/abs/2010.02502) (DDIM) by Jiaming Song, Chenlin Meng and Stefano Ermon.
The implementation is mostly based on the DDIM inversion definition of [Null-text Inversion for Editing Real Images using Guided Diffusion Models](https://arxiv.org/pdf/2211.09794.pdf)
## DDIMInverseScheduler
[[autodoc]] DDIMInverseScheduler

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specific language governing permissions and limitations under the License.
-->
# DDPMScheduler
# Denoising Diffusion Probabilistic Models (DDPM)
[Denoising Diffusion Probabilistic Models](https://huggingface.co/papers/2006.11239) (DDPM) by Jonathan Ho, Ajay Jain and Pieter Abbeel proposes a diffusion based model of the same name. In the context of the 🤗 Diffusers library, DDPM refers to the discrete denoising scheduler from the paper as well as the pipeline.
## Overview
The abstract from the paper is:
[Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2006.11239)
(DDPM) by Jonathan Ho, Ajay Jain and Pieter Abbeel proposes the diffusion based model of the same name, but in the context of the 🤗 Diffusers library, DDPM refers to the discrete denoising scheduler from the paper as well as the pipeline.
*We present high quality image synthesis results using diffusion probabilistic models, a class of latent variable models inspired by considerations from nonequilibrium thermodynamics. Our best results are obtained by training on a weighted variational bound designed according to a novel connection between diffusion probabilistic models and denoising score matching with Langevin dynamics, and our models naturally admit a progressive lossy decompression scheme that can be interpreted as a generalization of autoregressive decoding. On the unconditional CIFAR10 dataset, we obtain an Inception score of 9.46 and a state-of-the-art FID score of 3.17. On 256x256 LSUN, we obtain sample quality similar to ProgressiveGAN.*
The abstract of the paper is the following:
We present high quality image synthesis results using diffusion probabilistic models, a class of latent variable models inspired by considerations from nonequilibrium thermodynamics. Our best results are obtained by training on a weighted variational bound designed according to a novel connection between diffusion probabilistic models and denoising score matching with Langevin dynamics, and our models naturally admit a progressive lossy decompression scheme that can be interpreted as a generalization of autoregressive decoding. On the unconditional CIFAR10 dataset, we obtain an Inception score of 9.46 and a state-of-the-art FID score of 3.17. On 256x256 LSUN, we obtain sample quality similar to ProgressiveGAN.
The original paper can be found [here](https://arxiv.org/abs/2010.02502).
## DDPMScheduler
[[autodoc]] DDPMScheduler
## DDPMSchedulerOutput
[[autodoc]] schedulers.scheduling_ddpm.DDPMSchedulerOutput

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specific language governing permissions and limitations under the License.
-->
# DEISMultistepScheduler
# DEIS
Diffusion Exponential Integrator Sampler (DEIS) is proposed in [Fast Sampling of Diffusion Models with Exponential Integrator](https://huggingface.co/papers/2204.13902) by Qinsheng Zhang and Yongxin Chen. `DEISMultistepScheduler` is a fast high order solver for diffusion ordinary differential equations (ODEs).
Fast Sampling of Diffusion Models with Exponential Integrator.
This implementation modifies the polynomial fitting formula in log-rho space instead of the original linear `t` space in the DEIS paper. The modification enjoys closed-form coefficients for exponential multistep update instead of replying on the numerical solver.
## Overview
The abstract from the paper is:
*The past few years have witnessed the great success of Diffusion models~(DMs) in generating high-fidelity samples in generative modeling tasks. A major limitation of the DM is its notoriously slow sampling procedure which normally requires hundreds to thousands of time discretization steps of the learned diffusion process to reach the desired accuracy. Our goal is to develop a fast sampling method for DMs with a much less number of steps while retaining high sample quality. To this end, we systematically analyze the sampling procedure in DMs and identify key factors that affect the sample quality, among which the method of discretization is most crucial. By carefully examining the learned diffusion process, we propose Diffusion Exponential Integrator Sampler~(DEIS). It is based on the Exponential Integrator designed for discretizing ordinary differential equations (ODEs) and leverages a semilinear structure of the learned diffusion process to reduce the discretization error. The proposed method can be applied to any DMs and can generate high-fidelity samples in as few as 10 steps. In our experiments, it takes about 3 minutes on one A6000 GPU to generate 50k images from CIFAR10. Moreover, by directly using pre-trained DMs, we achieve the state-of-art sampling performance when the number of score function evaluation~(NFE) is limited, e.g., 4.17 FID with 10 NFEs, 3.37 FID, and 9.74 IS with only 15 NFEs on CIFAR10. Code is available at [this https URL](https://github.com/qsh-zh/deis).*
The original codebase can be found at [qsh-zh/deis](https://github.com/qsh-zh/deis).
## Tips
It is recommended to set `solver_order` to 2 or 3, while `solver_order=1` is equivalent to [`DDIMScheduler`].
Dynamic thresholding from [Imagen](https://huggingface.co/papers/2205.11487) is supported, and for pixel-space
diffusion models, you can set `thresholding=True` to use the dynamic thresholding.
Original paper can be found [here](https://arxiv.org/abs/2204.13902). The original implementation can be found [here](https://github.com/qsh-zh/deis).
## DEISMultistepScheduler
[[autodoc]] DEISMultistepScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput

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specific language governing permissions and limitations under the License.
-->
# KDPM2DiscreteScheduler
# DPM Discrete Scheduler inspired by Karras et. al paper
The `KDPM2DiscreteScheduler` is inspired by the [Elucidating the Design Space of Diffusion-Based Generative Models](https://huggingface.co/papers/2206.00364) paper, and the scheduler is ported from and created by [Katherine Crowson](https://github.com/crowsonkb/).
## Overview
The original codebase can be found at [crowsonkb/k-diffusion](https://github.com/crowsonkb/k-diffusion).
Inspired by [Karras et. al](https://arxiv.org/abs/2206.00364). Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
## KDPM2DiscreteScheduler
[[autodoc]] KDPM2DiscreteScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
[[autodoc]] KDPM2DiscreteScheduler

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specific language governing permissions and limitations under the License.
-->
# KDPM2AncestralDiscreteScheduler
# DPM Discrete Scheduler with ancestral sampling inspired by Karras et. al paper
The `KDPM2DiscreteScheduler` with ancestral sampling is inspired by the [Elucidating the Design Space of Diffusion-Based Generative Models](https://huggingface.co/papers/2206.00364) paper, and the scheduler is ported from and created by [Katherine Crowson](https://github.com/crowsonkb/).
## Overview
The original codebase can be found at [crowsonkb/k-diffusion](https://github.com/crowsonkb/k-diffusion).
Inspired by [Karras et. al](https://arxiv.org/abs/2206.00364). Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
## KDPM2AncestralDiscreteScheduler
[[autodoc]] KDPM2AncestralDiscreteScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
[[autodoc]] KDPM2AncestralDiscreteScheduler

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specific language governing permissions and limitations under the License.
-->
# DPMSolverSDEScheduler
# DPM Stochastic Scheduler inspired by Karras et. al paper
The `DPMSolverSDEScheduler` is inspired by the stochastic sampler from the [Elucidating the Design Space of Diffusion-Based Generative Models](https://huggingface.co/papers/2206.00364) paper, and the scheduler is ported from and created by [Katherine Crowson](https://github.com/crowsonkb/).
## Overview
Inspired by Stochastic Sampler from [Karras et. al](https://arxiv.org/abs/2206.00364).
Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
## DPMSolverSDEScheduler
[[autodoc]] DPMSolverSDEScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
[[autodoc]] DPMSolverSDEScheduler

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specific language governing permissions and limitations under the License.
-->
# EulerDiscreteScheduler
# Euler scheduler
The Euler scheduler (Algorithm 2) is from the [Elucidating the Design Space of Diffusion-Based Generative Models](https://huggingface.co/papers/2206.00364) paper by Karras et al. This is a fast scheduler which can often generate good outputs in 20-30 steps. The scheduler is based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L51) implementation by [Katherine Crowson](https://github.com/crowsonkb/).
## Overview
Euler scheduler (Algorithm 2) from the paper [Elucidating the Design Space of Diffusion-Based Generative Models](https://arxiv.org/abs/2206.00364) by Karras et al. (2022). Based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L51) implementation by Katherine Crowson.
Fast scheduler which often times generates good outputs with 20-30 steps.
## EulerDiscreteScheduler
[[autodoc]] EulerDiscreteScheduler
## EulerDiscreteSchedulerOutput
[[autodoc]] schedulers.scheduling_euler_discrete.EulerDiscreteSchedulerOutput
[[autodoc]] EulerDiscreteScheduler

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specific language governing permissions and limitations under the License.
-->
# EulerAncestralDiscreteScheduler
# Euler Ancestral scheduler
A scheduler that uses ancestral sampling with Euler method steps. This is a fast scheduler which can often generate good outputs in 20-30 steps. The scheduler is based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72) implementation by [Katherine Crowson](https://github.com/crowsonkb/).
## Overview
Ancestral sampling with Euler method steps. Based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72) implementation by Katherine Crowson.
Fast scheduler which often times generates good outputs with 20-30 steps.
## EulerAncestralDiscreteScheduler
[[autodoc]] EulerAncestralDiscreteScheduler
## EulerAncestralDiscreteSchedulerOutput
[[autodoc]] schedulers.scheduling_euler_ancestral_discrete.EulerAncestralDiscreteSchedulerOutput

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specific language governing permissions and limitations under the License.
-->
# HeunDiscreteScheduler
# Heun scheduler inspired by Karras et. al paper
The Heun scheduler (Algorithm 1) is from the [Elucidating the Design Space of Diffusion-Based Generative Models](https://huggingface.co/papers/2206.00364) paper by Karras et al. The scheduler is ported from the [k-diffusion](https://github.com/crowsonkb/k-diffusion) library and created by [Katherine Crowson](https://github.com/crowsonkb/).
## Overview
Algorithm 1 of [Karras et. al](https://arxiv.org/abs/2206.00364).
Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
## HeunDiscreteScheduler
[[autodoc]] HeunDiscreteScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
[[autodoc]] HeunDiscreteScheduler

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specific language governing permissions and limitations under the License.
-->
# IPNDMScheduler
# improved pseudo numerical methods for diffusion models (iPNDM)
`IPNDMScheduler` is a fourth-order Improved Pseudo Linear Multistep scheduler. The original implementation can be found at [crowsonkb/v-diffusion-pytorch](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296).
## Overview
Original implementation can be found [here](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296).
## IPNDMScheduler
[[autodoc]] IPNDMScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
[[autodoc]] IPNDMScheduler

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# Latent Consistency Model Multistep Scheduler
## Overview
Multistep and onestep scheduler (Algorithm 3) introduced alongside latent consistency models in the paper [Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference](https://arxiv.org/abs/2310.04378) by Simian Luo, Yiqin Tan, Longbo Huang, Jian Li, and Hang Zhao.
This scheduler should be able to generate good samples from [`LatentConsistencyModelPipeline`] in 1-8 steps.
## LCMScheduler
[[autodoc]] LCMScheduler

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specific language governing permissions and limitations under the License.
-->
# LMSDiscreteScheduler
# Linear multistep scheduler for discrete beta schedules
`LMSDiscreteScheduler` is a linear multistep scheduler for discrete beta schedules. The scheduler is ported from and created by [Katherine Crowson](https://github.com/crowsonkb/), and the original implementation can be found at [crowsonkb/k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181).
## Overview
Original implementation can be found [here](https://arxiv.org/abs/2206.00364).
## LMSDiscreteScheduler
[[autodoc]] LMSDiscreteScheduler
## LMSDiscreteSchedulerOutput
[[autodoc]] schedulers.scheduling_lms_discrete.LMSDiscreteSchedulerOutput
[[autodoc]] LMSDiscreteScheduler

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specific language governing permissions and limitations under the License.
-->
# DPMSolverMultistepScheduler
# Multistep DPM-Solver
`DPMSolverMultistep` is a multistep scheduler from [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095) by Cheng Lu, Yuhao Zhou, Fan Bao, Jianfei Chen, Chongxuan Li, and Jun Zhu.
## Overview
DPMSolver (and the improved version DPMSolver++) is a fast dedicated high-order solver for diffusion ODEs with convergence order guarantee. Empirically, DPMSolver sampling with only 20 steps can generate high-quality
samples, and it can generate quite good samples even in 10 steps.
## Tips
It is recommended to set `solver_order` to 2 for guide sampling, and `solver_order=3` for unconditional sampling.
Dynamic thresholding from Imagen (https://huggingface.co/papers/2205.11487) is supported, and for pixel-space
diffusion models, you can set both `algorithm_type="dpmsolver++"` and `thresholding=True` to use the dynamic
thresholding. This thresholding method is unsuitable for latent-space diffusion models such as
Stable Diffusion.
The SDE variant of DPMSolver and DPM-Solver++ is also supported, but only for the first and second-order solvers. This is a fast SDE solver for the reverse diffusion SDE. It is recommended to use the second-order `sde-dpmsolver++`.
Original paper can be found [here](https://arxiv.org/abs/2206.00927) and the [improved version](https://arxiv.org/abs/2211.01095). The original implementation can be found [here](https://github.com/LuChengTHU/dpm-solver).
## DPMSolverMultistepScheduler
[[autodoc]] DPMSolverMultistepScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
[[autodoc]] DPMSolverMultistepScheduler

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specific language governing permissions and limitations under the License.
-->
# DPMSolverMultistepInverse
# Inverse Multistep DPM-Solver (DPMSolverMultistepInverse)
`DPMSolverMultistepInverse` is the inverted scheduler from [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095) by Cheng Lu, Yuhao Zhou, Fan Bao, Jianfei Chen, Chongxuan Li, and Jun Zhu.
## Overview
The implementation is mostly based on the DDIM inversion definition of [Null-text Inversion for Editing Real Images using Guided Diffusion Models](https://huggingface.co/papers/2211.09794.pdf) and notebook implementation of the [`DiffEdit`] latent inversion from [Xiang-cd/DiffEdit-stable-diffusion](https://github.com/Xiang-cd/DiffEdit-stable-diffusion/blob/main/diffedit.ipynb).
## Tips
Dynamic thresholding from Imagen (https://huggingface.co/papers/2205.11487) is supported, and for pixel-space
diffusion models, you can set both `algorithm_type="dpmsolver++"` and `thresholding=True` to use the dynamic
thresholding. This thresholding method is unsuitable for latent-space diffusion models such as
Stable Diffusion.
This scheduler is the inverted scheduler of [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://arxiv.org/abs/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models
](https://arxiv.org/abs/2211.01095) by Cheng Lu, Yuhao Zhou, Fan Bao, Jianfei Chen, Chongxuan Li, and Jun Zhu.
The implementation is mostly based on the DDIM inversion definition of [Null-text Inversion for Editing Real Images using Guided Diffusion Models](https://arxiv.org/pdf/2211.09794.pdf) and the ad-hoc notebook implementation for DiffEdit latent inversion [here](https://github.com/Xiang-cd/DiffEdit-stable-diffusion/blob/main/diffedit.ipynb).
## DPMSolverMultistepInverseScheduler
[[autodoc]] DPMSolverMultistepInverseScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput

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# Schedulers
🤗 Diffusers provides many scheduler functions for the diffusion process. A scheduler takes a model's output (the sample which the diffusion process is iterating on) and a timestep to return a denoised sample. The timestep is important because it dictates where in the diffusion process the step is; data is generated by iterating forward *n* timesteps and inference occurs by propagating backward through the timesteps. Based on the timestep, a scheduler may be *discrete* in which case the timestep is an `int` or *continuous* in which case the timestep is a `float`.
Diffusers contains multiple pre-built schedule functions for the diffusion process.
Depending on the context, a scheduler defines how to iteratively add noise to an image or how to update a sample based on a model's output:
## What is a scheduler?
- during *training*, a scheduler adds noise (there are different algorithms for how to add noise) to a sample to train a diffusion model
- during *inference*, a scheduler defines how to update a sample based on a pretrained model's output
The schedule functions, denoted *Schedulers* in the library take in the output of a trained model, a sample which the diffusion process is iterating on, and a timestep to return a denoised sample. That's why schedulers may also be called *Samplers* in other diffusion models implementations.
Many schedulers are implemented from the [k-diffusion](https://github.com/crowsonkb/k-diffusion) library by [Katherine Crowson](https://github.com/crowsonkb/), and they're also widely used in A1111. To help you map the schedulers from k-diffusion and A1111 to the schedulers in 🤗 Diffusers, take a look at the table below:
- Schedulers define the methodology for iteratively adding noise to an image or for updating a sample based on model outputs.
- adding noise in different manners represent the algorithmic processes to train a diffusion model by adding noise to images.
- for inference, the scheduler defines how to update a sample based on an output from a pretrained model.
- Schedulers are often defined by a *noise schedule* and an *update rule* to solve the differential equation solution.
| A1111/k-diffusion | 🤗 Diffusers | Usage |
|---------------------|-------------------------------------|---------------------------------------------------------------------------------------------------------------|
| DPM++ 2M | [`DPMSolverMultistepScheduler`] | |
| DPM++ 2M Karras | [`DPMSolverMultistepScheduler`] | init with `use_karras_sigmas=True` |
| DPM++ 2M SDE | [`DPMSolverMultistepScheduler`] | init with `algorithm_type="sde-dpmsolver++"` |
| DPM++ 2M SDE Karras | [`DPMSolverMultistepScheduler`] | init with `use_karras_sigmas=True` and `algorithm_type="sde-dpmsolver++"` |
| DPM++ 2S a | N/A | very similar to `DPMSolverSinglestepScheduler` |
| DPM++ 2S a Karras | N/A | very similar to `DPMSolverSinglestepScheduler(use_karras_sigmas=True, ...)` |
| DPM++ SDE | [`DPMSolverSinglestepScheduler`] | |
| DPM++ SDE Karras | [`DPMSolverSinglestepScheduler`] | init with `use_karras_sigmas=True` |
| DPM2 | [`KDPM2DiscreteScheduler`] | |
| DPM2 Karras | [`KDPM2DiscreteScheduler`] | init with `use_karras_sigmas=True` |
| DPM2 a | [`KDPM2AncestralDiscreteScheduler`] | |
| DPM2 a Karras | [`KDPM2AncestralDiscreteScheduler`] | init with `use_karras_sigmas=True` |
| DPM adaptive | N/A | |
| DPM fast | N/A | |
| Euler | [`EulerDiscreteScheduler`] | |
| Euler a | [`EulerAncestralDiscreteScheduler`] | |
| Heun | [`HeunDiscreteScheduler`] | |
| LMS | [`LMSDiscreteScheduler`] | |
| LMS Karras | [`LMSDiscreteScheduler`] | init with `use_karras_sigmas=True` |
| N/A | [`DEISMultistepScheduler`] | |
| N/A | [`UniPCMultistepScheduler`] | |
### Discrete versus continuous schedulers
All schedulers are built from the base [`SchedulerMixin`] class which implements low level utilities shared by all schedulers.
All schedulers take in a timestep to predict the updated version of the sample being diffused.
The timesteps dictate where in the diffusion process the step is, where data is generated by iterating forward in time and inference is executed by propagating backwards through timesteps.
Different algorithms use timesteps that can be discrete (accepting `int` inputs), such as the [`DDPMScheduler`] or [`PNDMScheduler`], or continuous (accepting `float` inputs), such as the score-based schedulers [`ScoreSdeVeScheduler`] or [`ScoreSdeVpScheduler`].
## SchedulerMixin
## Designing Re-usable schedulers
The core design principle between the schedule functions is to be model, system, and framework independent.
This allows for rapid experimentation and cleaner abstractions in the code, where the model prediction is separated from the sample update.
To this end, the design of schedulers is such that:
- Schedulers can be used interchangeably between diffusion models in inference to find the preferred trade-off between speed and generation quality.
- Schedulers are currently by default in PyTorch, but are designed to be framework independent (partial Jax support currently exists).
- Many diffusion pipelines, such as [`StableDiffusionPipeline`] and [`DiTPipeline`] can use any of [`KarrasDiffusionSchedulers`]
## Schedulers Summary
The following table summarizes all officially supported schedulers, their corresponding paper
| Scheduler | Paper |
|---|---|
| [ddim](./ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) |
| [ddim_inverse](./ddim_inverse) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) |
| [ddpm](./ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) |
| [deis](./deis) | [**DEISMultistepScheduler**](https://arxiv.org/abs/2204.13902) |
| [singlestep_dpm_solver](./singlestep_dpm_solver) | [**Singlestep DPM-Solver**](https://arxiv.org/abs/2206.00927) |
| [multistep_dpm_solver](./multistep_dpm_solver) | [**Multistep DPM-Solver**](https://arxiv.org/abs/2206.00927) |
| [heun](./heun) | [**Heun scheduler inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
| [dpm_discrete](./dpm_discrete) | [**DPM Discrete Scheduler inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
| [dpm_discrete_ancestral](./dpm_discrete_ancestral) | [**DPM Discrete Scheduler with ancestral sampling inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
| [stochastic_karras_ve](./stochastic_karras_ve) | [**Variance exploding, stochastic sampling from Karras et. al**](https://arxiv.org/abs/2206.00364) |
| [lms_discrete](./lms_discrete) | [**Linear multistep scheduler for discrete beta schedules**](https://arxiv.org/abs/2206.00364) |
| [pndm](./pndm) | [**Pseudo numerical methods for diffusion models (PNDM)**](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181) |
| [score_sde_ve](./score_sde_ve) | [**variance exploding stochastic differential equation (VE-SDE) scheduler**](https://arxiv.org/abs/2011.13456) |
| [ipndm](./ipndm) | [**improved pseudo numerical methods for diffusion models (iPNDM)**](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296) |
| [score_sde_vp](./score_sde_vp) | [**Variance preserving stochastic differential equation (VP-SDE) scheduler**](https://arxiv.org/abs/2011.13456) |
| [euler](./euler) | [**Euler scheduler**](https://arxiv.org/abs/2206.00364) |
| [euler_ancestral](./euler_ancestral) | [**Euler Ancestral scheduler**](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72) |
| [vq_diffusion](./vq_diffusion) | [**VQDiffusionScheduler**](https://arxiv.org/abs/2111.14822) |
| [unipc](./unipc) | [**UniPCMultistepScheduler**](https://arxiv.org/abs/2302.04867) |
| [repaint](./repaint) | [**RePaint scheduler**](https://arxiv.org/abs/2201.09865) |
## API
The core API for any new scheduler must follow a limited structure.
- Schedulers should provide one or more `def step(...)` functions that should be called to update the generated sample iteratively.
- Schedulers should provide a `set_timesteps(...)` method that configures the parameters of a schedule function for a specific inference task.
- Schedulers should be framework-specific.
The base class [`SchedulerMixin`] implements low level utilities used by multiple schedulers.
### SchedulerMixin
[[autodoc]] SchedulerMixin
## SchedulerOutput
### SchedulerOutput
The class [`SchedulerOutput`] contains the outputs from any schedulers `step(...)` call.
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
## KarrasDiffusionSchedulers
### KarrasDiffusionSchedulers
[`KarrasDiffusionSchedulers`] are a broad generalization of schedulers in 🤗 Diffusers. The schedulers in this class are distinguished at a high level by their noise sampling strategy, the type of network and scaling, the training strategy, and how the loss is weighed.
`KarrasDiffusionSchedulers` encompasses the main generalization of schedulers in Diffusers. The schedulers in this class are distinguished, at a high level, by their noise sampling strategy; the type of network and scaling; and finally the training strategy or how the loss is weighed.
The different schedulers in this class, depending on the ordinary differential equations (ODE) solver type, fall into the above taxonomy and provide a good abstraction for the design of the main schedulers implemented in 🤗 Diffusers. The schedulers in this class are given [here](https://github.com/huggingface/diffusers/blob/a69754bb879ed55b9b6dc9dd0b3cf4fa4124c765/src/diffusers/schedulers/scheduling_utils.py#L32).
The different schedulers, depending on the type of ODE solver, fall into the above taxonomy and provide a good abstraction for the design of the main schedulers implemented in Diffusers. The schedulers in this class are given below:
## PushToHubMixin
[[autodoc]] utils.PushToHubMixin
[[autodoc]] schedulers.scheduling_utils.KarrasDiffusionSchedulers

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