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v0.13.1-pa
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add_schedu
| Author | SHA1 | Date | |
|---|---|---|---|
|
|
b444122ab0 |
17
.github/ISSUE_TEMPLATE/bug-report.yml
vendored
17
.github/ISSUE_TEMPLATE/bug-report.yml
vendored
@@ -5,20 +5,7 @@ body:
|
||||
- type: markdown
|
||||
attributes:
|
||||
value: |
|
||||
Thanks a lot for taking the time to file this issue 🤗.
|
||||
Issues do not only help to improve the library, but also publicly document common problems, questions, workflows for the whole community!
|
||||
Thus, issues are of the same importance as pull requests when contributing to this library ❤️.
|
||||
In order to make your issue as **useful for the community as possible**, let's try to stick to some simple guidelines:
|
||||
- 1. Please try to be as precise and concise as possible.
|
||||
*Give your issue a fitting title. Assume that someone which very limited knowledge of diffusers can understand your issue. Add links to the source code, documentation other issues, pull requests etc...*
|
||||
- 2. If your issue is about something not working, **always** provide a reproducible code snippet. The reader should be able to reproduce your issue by **only copy-pasting your code snippet into a Python shell**.
|
||||
*The community cannot solve your issue if it cannot reproduce it. If your bug is related to training, add your training script and make everything needed to train public. Otherwise, just add a simple Python code snippet.*
|
||||
- 3. Add the **minimum amount of code / context that is needed to understand, reproduce your issue**.
|
||||
*Make the life of maintainers easy. `diffusers` is getting many issues every day. Make sure your issue is about one bug and one bug only. Make sure you add only the context, code needed to understand your issues - nothing more. Generally, every issue is a way of documenting this library, try to make it a good documentation entry.*
|
||||
- type: markdown
|
||||
attributes:
|
||||
value: |
|
||||
For more in-detail information on how to write good issues you can have a look [here](https://huggingface.co/course/chapter8/5?fw=pt)
|
||||
Thanks for taking the time to fill out this bug report!
|
||||
- type: textarea
|
||||
id: bug-description
|
||||
attributes:
|
||||
@@ -33,8 +20,6 @@ body:
|
||||
label: Reproduction
|
||||
description: Please provide a minimal reproducible code which we can copy/paste and reproduce the issue.
|
||||
placeholder: Reproduction
|
||||
validations:
|
||||
required: true
|
||||
- type: textarea
|
||||
id: logs
|
||||
attributes:
|
||||
|
||||
8
.github/ISSUE_TEMPLATE/config.yml
vendored
8
.github/ISSUE_TEMPLATE/config.yml
vendored
@@ -1,7 +1,7 @@
|
||||
contact_links:
|
||||
- name: Forum
|
||||
url: https://discuss.huggingface.co/c/discussion-related-to-httpsgithubcomhuggingfacediffusers/63
|
||||
about: General usage questions and community discussions
|
||||
- name: Blank issue
|
||||
url: https://github.com/huggingface/diffusers/issues/new
|
||||
about: Other
|
||||
- name: Forum
|
||||
url: https://discuss.huggingface.co/
|
||||
about: General usage questions and community discussions
|
||||
about: Please note that the Forum is in most places the right place for discussions
|
||||
|
||||
50
.github/workflows/build_docker_images.yml
vendored
50
.github/workflows/build_docker_images.yml
vendored
@@ -1,50 +0,0 @@
|
||||
name: Build Docker images (nightly)
|
||||
|
||||
on:
|
||||
workflow_dispatch:
|
||||
schedule:
|
||||
- cron: "0 0 * * *" # every day at midnight
|
||||
|
||||
concurrency:
|
||||
group: docker-image-builds
|
||||
cancel-in-progress: false
|
||||
|
||||
env:
|
||||
REGISTRY: diffusers
|
||||
|
||||
jobs:
|
||||
build-docker-images:
|
||||
runs-on: ubuntu-latest
|
||||
|
||||
permissions:
|
||||
contents: read
|
||||
packages: write
|
||||
|
||||
strategy:
|
||||
fail-fast: false
|
||||
matrix:
|
||||
image-name:
|
||||
- diffusers-pytorch-cpu
|
||||
- diffusers-pytorch-cuda
|
||||
- diffusers-flax-cpu
|
||||
- diffusers-flax-tpu
|
||||
- diffusers-onnxruntime-cpu
|
||||
- diffusers-onnxruntime-cuda
|
||||
|
||||
steps:
|
||||
- name: Checkout repository
|
||||
uses: actions/checkout@v3
|
||||
|
||||
- name: Login to Docker Hub
|
||||
uses: docker/login-action@v2
|
||||
with:
|
||||
username: ${{ env.REGISTRY }}
|
||||
password: ${{ secrets.DOCKERHUB_TOKEN }}
|
||||
|
||||
- name: Build and push
|
||||
uses: docker/build-push-action@v3
|
||||
with:
|
||||
no-cache: true
|
||||
context: ./docker/${{ matrix.image-name }}
|
||||
push: true
|
||||
tags: ${{ env.REGISTRY }}/${{ matrix.image-name }}:latest
|
||||
1
.github/workflows/build_documentation.yml
vendored
1
.github/workflows/build_documentation.yml
vendored
@@ -13,6 +13,5 @@ jobs:
|
||||
with:
|
||||
commit_sha: ${{ github.sha }}
|
||||
package: diffusers
|
||||
languages: en ko
|
||||
secrets:
|
||||
token: ${{ secrets.HUGGINGFACE_PUSH }}
|
||||
|
||||
1
.github/workflows/build_pr_documentation.yml
vendored
1
.github/workflows/build_pr_documentation.yml
vendored
@@ -14,4 +14,3 @@ jobs:
|
||||
commit_sha: ${{ github.event.pull_request.head.sha }}
|
||||
pr_number: ${{ github.event.number }}
|
||||
package: diffusers
|
||||
languages: en ko
|
||||
|
||||
162
.github/workflows/nightly_tests.yml
vendored
162
.github/workflows/nightly_tests.yml
vendored
@@ -1,162 +0,0 @@
|
||||
name: Nightly tests on main
|
||||
|
||||
on:
|
||||
schedule:
|
||||
- cron: "0 0 * * *" # every day at midnight
|
||||
|
||||
env:
|
||||
DIFFUSERS_IS_CI: yes
|
||||
HF_HOME: /mnt/cache
|
||||
OMP_NUM_THREADS: 8
|
||||
MKL_NUM_THREADS: 8
|
||||
PYTEST_TIMEOUT: 600
|
||||
RUN_SLOW: yes
|
||||
RUN_NIGHTLY: yes
|
||||
|
||||
jobs:
|
||||
run_nightly_tests:
|
||||
strategy:
|
||||
fail-fast: false
|
||||
matrix:
|
||||
config:
|
||||
- name: Nightly PyTorch CUDA tests on Ubuntu
|
||||
framework: pytorch
|
||||
runner: docker-gpu
|
||||
image: diffusers/diffusers-pytorch-cuda
|
||||
report: torch_cuda
|
||||
- name: Nightly Flax TPU tests on Ubuntu
|
||||
framework: flax
|
||||
runner: docker-tpu
|
||||
image: diffusers/diffusers-flax-tpu
|
||||
report: flax_tpu
|
||||
- name: Nightly ONNXRuntime CUDA tests on Ubuntu
|
||||
framework: onnxruntime
|
||||
runner: docker-gpu
|
||||
image: diffusers/diffusers-onnxruntime-cuda
|
||||
report: onnx_cuda
|
||||
|
||||
name: ${{ matrix.config.name }}
|
||||
|
||||
runs-on: ${{ matrix.config.runner }}
|
||||
|
||||
container:
|
||||
image: ${{ matrix.config.image }}
|
||||
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ ${{ matrix.config.runner == 'docker-tpu' && '--privileged' || '--gpus 0'}}
|
||||
|
||||
defaults:
|
||||
run:
|
||||
shell: bash
|
||||
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
uses: actions/checkout@v3
|
||||
with:
|
||||
fetch-depth: 2
|
||||
|
||||
- name: NVIDIA-SMI
|
||||
if: ${{ matrix.config.runner == 'docker-gpu' }}
|
||||
run: |
|
||||
nvidia-smi
|
||||
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
python -m pip install -e .[quality,test]
|
||||
python -m pip install -U git+https://github.com/huggingface/transformers
|
||||
python -m pip install git+https://github.com/huggingface/accelerate
|
||||
|
||||
- name: Environment
|
||||
run: |
|
||||
python utils/print_env.py
|
||||
|
||||
- name: Run nightly PyTorch CUDA tests
|
||||
if: ${{ matrix.config.framework == 'pytorch' }}
|
||||
env:
|
||||
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "not Flax and not Onnx" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Run nightly Flax TPU tests
|
||||
if: ${{ matrix.config.framework == 'flax' }}
|
||||
env:
|
||||
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 0 \
|
||||
-s -v -k "Flax" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Run nightly ONNXRuntime CUDA tests
|
||||
if: ${{ matrix.config.framework == 'onnxruntime' }}
|
||||
env:
|
||||
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "Onnx" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
|
||||
|
||||
- name: Test suite reports artifacts
|
||||
if: ${{ always() }}
|
||||
uses: actions/upload-artifact@v2
|
||||
with:
|
||||
name: ${{ matrix.config.report }}_test_reports
|
||||
path: reports
|
||||
|
||||
run_nightly_tests_apple_m1:
|
||||
name: Nightly PyTorch MPS tests on MacOS
|
||||
runs-on: [ self-hosted, apple-m1 ]
|
||||
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
uses: actions/checkout@v3
|
||||
with:
|
||||
fetch-depth: 2
|
||||
|
||||
- name: Clean checkout
|
||||
shell: arch -arch arm64 bash {0}
|
||||
run: |
|
||||
git clean -fxd
|
||||
|
||||
- name: Setup miniconda
|
||||
uses: ./.github/actions/setup-miniconda
|
||||
with:
|
||||
python-version: 3.9
|
||||
|
||||
- name: Install dependencies
|
||||
shell: arch -arch arm64 bash {0}
|
||||
run: |
|
||||
${CONDA_RUN} python -m pip install --upgrade pip
|
||||
${CONDA_RUN} python -m pip install -e .[quality,test]
|
||||
${CONDA_RUN} python -m pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
|
||||
${CONDA_RUN} python -m pip install git+https://github.com/huggingface/accelerate
|
||||
|
||||
- name: Environment
|
||||
shell: arch -arch arm64 bash {0}
|
||||
run: |
|
||||
${CONDA_RUN} python utils/print_env.py
|
||||
|
||||
- name: Run nightly PyTorch tests on M1 (MPS)
|
||||
shell: arch -arch arm64 bash {0}
|
||||
env:
|
||||
HF_HOME: /System/Volumes/Data/mnt/cache
|
||||
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
|
||||
run: |
|
||||
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps tests/
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
run: cat reports/tests_torch_mps_failures_short.txt
|
||||
|
||||
- name: Test suite reports artifacts
|
||||
if: ${{ always() }}
|
||||
uses: actions/upload-artifact@v2
|
||||
with:
|
||||
name: torch_mps_test_reports
|
||||
path: reports
|
||||
5
.github/workflows/pr_quality.yml
vendored
5
.github/workflows/pr_quality.yml
vendored
@@ -27,8 +27,9 @@ jobs:
|
||||
pip install .[quality]
|
||||
- name: Check quality
|
||||
run: |
|
||||
black --check examples tests src utils scripts
|
||||
ruff examples tests src utils scripts
|
||||
black --check --preview examples tests src utils scripts
|
||||
isort --check-only examples tests src utils scripts
|
||||
flake8 examples tests src utils scripts
|
||||
doc-builder style src/diffusers docs/source --max_len 119 --check_only --path_to_docs docs/source
|
||||
|
||||
check_repository_consistency:
|
||||
|
||||
100
.github/workflows/pr_tests.yml
vendored
100
.github/workflows/pr_tests.yml
vendored
@@ -1,4 +1,4 @@
|
||||
name: Fast tests for PRs
|
||||
name: Run fast tests
|
||||
|
||||
on:
|
||||
pull_request:
|
||||
@@ -10,50 +10,19 @@ concurrency:
|
||||
cancel-in-progress: true
|
||||
|
||||
env:
|
||||
DIFFUSERS_IS_CI: yes
|
||||
OMP_NUM_THREADS: 4
|
||||
MKL_NUM_THREADS: 4
|
||||
OMP_NUM_THREADS: 8
|
||||
MKL_NUM_THREADS: 8
|
||||
PYTEST_TIMEOUT: 60
|
||||
MPS_TORCH_VERSION: 1.13.0
|
||||
|
||||
jobs:
|
||||
run_fast_tests:
|
||||
strategy:
|
||||
fail-fast: false
|
||||
matrix:
|
||||
config:
|
||||
- name: Fast PyTorch CPU tests on Ubuntu
|
||||
framework: pytorch
|
||||
runner: docker-cpu
|
||||
image: diffusers/diffusers-pytorch-cpu
|
||||
report: torch_cpu
|
||||
- name: Fast Flax CPU tests on Ubuntu
|
||||
framework: flax
|
||||
runner: docker-cpu
|
||||
image: diffusers/diffusers-flax-cpu
|
||||
report: flax_cpu
|
||||
- name: Fast ONNXRuntime CPU tests on Ubuntu
|
||||
framework: onnxruntime
|
||||
runner: docker-cpu
|
||||
image: diffusers/diffusers-onnxruntime-cpu
|
||||
report: onnx_cpu
|
||||
- name: PyTorch Example CPU tests on Ubuntu
|
||||
framework: pytorch_examples
|
||||
runner: docker-cpu
|
||||
image: diffusers/diffusers-pytorch-cpu
|
||||
report: torch_cpu
|
||||
|
||||
name: ${{ matrix.config.name }}
|
||||
|
||||
runs-on: ${{ matrix.config.runner }}
|
||||
|
||||
run_tests_cpu:
|
||||
name: CPU tests on Ubuntu
|
||||
runs-on: [ self-hosted, docker-gpu ]
|
||||
container:
|
||||
image: ${{ matrix.config.image }}
|
||||
image: python:3.7
|
||||
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
|
||||
|
||||
defaults:
|
||||
run:
|
||||
shell: bash
|
||||
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
uses: actions/checkout@v3
|
||||
@@ -62,59 +31,32 @@ jobs:
|
||||
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
apt-get update && apt-get install libsndfile1-dev -y
|
||||
python -m pip install --upgrade pip
|
||||
python -m pip install torch --extra-index-url https://download.pytorch.org/whl/cpu
|
||||
python -m pip install -e .[quality,test]
|
||||
python -m pip install -U git+https://github.com/huggingface/transformers
|
||||
python -m pip install git+https://github.com/huggingface/accelerate
|
||||
|
||||
- name: Environment
|
||||
run: |
|
||||
python utils/print_env.py
|
||||
|
||||
- name: Run fast PyTorch CPU tests
|
||||
if: ${{ matrix.config.framework == 'pytorch' }}
|
||||
- name: Run all fast tests on CPU
|
||||
run: |
|
||||
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "not Flax and not Onnx" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Run fast Flax TPU tests
|
||||
if: ${{ matrix.config.framework == 'flax' }}
|
||||
run: |
|
||||
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "Flax" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Run fast ONNXRuntime CPU tests
|
||||
if: ${{ matrix.config.framework == 'onnxruntime' }}
|
||||
run: |
|
||||
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "Onnx" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Run example PyTorch CPU tests
|
||||
if: ${{ matrix.config.framework == 'pytorch_examples' }}
|
||||
run: |
|
||||
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
examples/test_examples.py
|
||||
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=tests_torch_cpu tests/
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
|
||||
run: cat reports/tests_torch_cpu_failures_short.txt
|
||||
|
||||
- name: Test suite reports artifacts
|
||||
if: ${{ always() }}
|
||||
uses: actions/upload-artifact@v2
|
||||
with:
|
||||
name: pr_${{ matrix.config.report }}_test_reports
|
||||
name: pr_torch_cpu_test_reports
|
||||
path: reports
|
||||
|
||||
run_fast_tests_apple_m1:
|
||||
name: Fast PyTorch MPS tests on MacOS
|
||||
run_tests_apple_m1:
|
||||
name: MPS tests on Apple M1
|
||||
runs-on: [ self-hosted, apple-m1 ]
|
||||
|
||||
steps:
|
||||
@@ -138,22 +80,18 @@ jobs:
|
||||
run: |
|
||||
${CONDA_RUN} python -m pip install --upgrade pip
|
||||
${CONDA_RUN} python -m pip install -e .[quality,test]
|
||||
${CONDA_RUN} python -m pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
|
||||
${CONDA_RUN} python -m pip install --pre torch==${MPS_TORCH_VERSION} --extra-index-url https://download.pytorch.org/whl/test/cpu
|
||||
${CONDA_RUN} python -m pip install git+https://github.com/huggingface/accelerate
|
||||
${CONDA_RUN} python -m pip install -U git+https://github.com/huggingface/transformers
|
||||
|
||||
- name: Environment
|
||||
shell: arch -arch arm64 bash {0}
|
||||
run: |
|
||||
${CONDA_RUN} python utils/print_env.py
|
||||
|
||||
- name: Run fast PyTorch tests on M1 (MPS)
|
||||
- name: Run all fast tests on MPS
|
||||
shell: arch -arch arm64 bash {0}
|
||||
env:
|
||||
HF_HOME: /System/Volumes/Data/mnt/cache
|
||||
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
|
||||
run: |
|
||||
${CONDA_RUN} python -m pytest -n 0 -s -v --make-reports=tests_torch_mps tests/
|
||||
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps tests/
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
|
||||
98
.github/workflows/push_tests.yml
vendored
98
.github/workflows/push_tests.yml
vendored
@@ -1,4 +1,4 @@
|
||||
name: Slow tests on main
|
||||
name: Run all tests
|
||||
|
||||
on:
|
||||
push:
|
||||
@@ -6,46 +6,19 @@ on:
|
||||
- main
|
||||
|
||||
env:
|
||||
DIFFUSERS_IS_CI: yes
|
||||
HF_HOME: /mnt/cache
|
||||
OMP_NUM_THREADS: 8
|
||||
MKL_NUM_THREADS: 8
|
||||
PYTEST_TIMEOUT: 600
|
||||
PYTEST_TIMEOUT: 1000
|
||||
RUN_SLOW: yes
|
||||
|
||||
jobs:
|
||||
run_slow_tests:
|
||||
strategy:
|
||||
fail-fast: false
|
||||
matrix:
|
||||
config:
|
||||
- name: Slow PyTorch CUDA tests on Ubuntu
|
||||
framework: pytorch
|
||||
runner: docker-gpu
|
||||
image: diffusers/diffusers-pytorch-cuda
|
||||
report: torch_cuda
|
||||
- name: Slow Flax TPU tests on Ubuntu
|
||||
framework: flax
|
||||
runner: docker-tpu
|
||||
image: diffusers/diffusers-flax-tpu
|
||||
report: flax_tpu
|
||||
- name: Slow ONNXRuntime CUDA tests on Ubuntu
|
||||
framework: onnxruntime
|
||||
runner: docker-gpu
|
||||
image: diffusers/diffusers-onnxruntime-cuda
|
||||
report: onnx_cuda
|
||||
|
||||
name: ${{ matrix.config.name }}
|
||||
|
||||
runs-on: ${{ matrix.config.runner }}
|
||||
|
||||
run_tests_single_gpu:
|
||||
name: Diffusers tests
|
||||
runs-on: [ self-hosted, docker-gpu, single-gpu ]
|
||||
container:
|
||||
image: ${{ matrix.config.image }}
|
||||
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ ${{ matrix.config.runner == 'docker-tpu' && '--privileged' || '--gpus 0'}}
|
||||
|
||||
defaults:
|
||||
run:
|
||||
shell: bash
|
||||
image: nvcr.io/nvidia/pytorch:22.07-py3
|
||||
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache
|
||||
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
@@ -54,69 +27,44 @@ jobs:
|
||||
fetch-depth: 2
|
||||
|
||||
- name: NVIDIA-SMI
|
||||
if : ${{ matrix.config.runner == 'docker-gpu' }}
|
||||
run: |
|
||||
nvidia-smi
|
||||
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
python -m pip install --upgrade pip
|
||||
python -m pip uninstall -y torch torchvision torchtext
|
||||
python -m pip install torch --extra-index-url https://download.pytorch.org/whl/cu117
|
||||
python -m pip install -e .[quality,test]
|
||||
python -m pip install -U git+https://github.com/huggingface/transformers
|
||||
python -m pip install git+https://github.com/huggingface/accelerate
|
||||
|
||||
- name: Environment
|
||||
run: |
|
||||
python utils/print_env.py
|
||||
|
||||
- name: Run slow PyTorch CUDA tests
|
||||
if: ${{ matrix.config.framework == 'pytorch' }}
|
||||
- name: Run all (incl. slow) tests on GPU
|
||||
env:
|
||||
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "not Flax and not Onnx" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Run slow Flax TPU tests
|
||||
if: ${{ matrix.config.framework == 'flax' }}
|
||||
env:
|
||||
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 0 \
|
||||
-s -v -k "Flax" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Run slow ONNXRuntime CUDA tests
|
||||
if: ${{ matrix.config.framework == 'onnxruntime' }}
|
||||
env:
|
||||
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "Onnx" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=tests_torch_gpu tests/
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
|
||||
run: cat reports/tests_torch_gpu_failures_short.txt
|
||||
|
||||
- name: Test suite reports artifacts
|
||||
if: ${{ always() }}
|
||||
uses: actions/upload-artifact@v2
|
||||
with:
|
||||
name: ${{ matrix.config.report }}_test_reports
|
||||
name: torch_test_reports
|
||||
path: reports
|
||||
|
||||
run_examples_tests:
|
||||
name: Examples PyTorch CUDA tests on Ubuntu
|
||||
|
||||
runs-on: docker-gpu
|
||||
|
||||
run_examples_single_gpu:
|
||||
name: Examples tests
|
||||
runs-on: [ self-hosted, docker-gpu, single-gpu ]
|
||||
container:
|
||||
image: diffusers/diffusers-pytorch-cuda
|
||||
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
|
||||
image: nvcr.io/nvidia/pytorch:22.07-py3
|
||||
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache
|
||||
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
@@ -130,9 +78,11 @@ jobs:
|
||||
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
python -m pip install --upgrade pip
|
||||
python -m pip uninstall -y torch torchvision torchtext
|
||||
python -m pip install torch --extra-index-url https://download.pytorch.org/whl/cu117
|
||||
python -m pip install -e .[quality,test,training]
|
||||
python -m pip install git+https://github.com/huggingface/accelerate
|
||||
python -m pip install -U git+https://github.com/huggingface/transformers
|
||||
|
||||
- name: Environment
|
||||
run: |
|
||||
@@ -142,11 +92,11 @@ jobs:
|
||||
env:
|
||||
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=examples_torch_cuda examples/
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=examples_torch_gpu examples/
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
run: cat reports/examples_torch_cuda_failures_short.txt
|
||||
run: cat reports/examples_torch_gpu_failures_short.txt
|
||||
|
||||
- name: Test suite reports artifacts
|
||||
if: ${{ always() }}
|
||||
|
||||
165
.github/workflows/push_tests_fast.yml
vendored
165
.github/workflows/push_tests_fast.yml
vendored
@@ -1,165 +0,0 @@
|
||||
name: Slow tests on main
|
||||
|
||||
on:
|
||||
push:
|
||||
branches:
|
||||
- main
|
||||
|
||||
env:
|
||||
DIFFUSERS_IS_CI: yes
|
||||
HF_HOME: /mnt/cache
|
||||
OMP_NUM_THREADS: 8
|
||||
MKL_NUM_THREADS: 8
|
||||
PYTEST_TIMEOUT: 600
|
||||
RUN_SLOW: no
|
||||
|
||||
jobs:
|
||||
run_fast_tests:
|
||||
strategy:
|
||||
fail-fast: false
|
||||
matrix:
|
||||
config:
|
||||
- name: Fast PyTorch CPU tests on Ubuntu
|
||||
framework: pytorch
|
||||
runner: docker-cpu
|
||||
image: diffusers/diffusers-pytorch-cpu
|
||||
report: torch_cpu
|
||||
- name: Fast Flax CPU tests on Ubuntu
|
||||
framework: flax
|
||||
runner: docker-cpu
|
||||
image: diffusers/diffusers-flax-cpu
|
||||
report: flax_cpu
|
||||
- name: Fast ONNXRuntime CPU tests on Ubuntu
|
||||
framework: onnxruntime
|
||||
runner: docker-cpu
|
||||
image: diffusers/diffusers-onnxruntime-cpu
|
||||
report: onnx_cpu
|
||||
- name: PyTorch Example CPU tests on Ubuntu
|
||||
framework: pytorch_examples
|
||||
runner: docker-cpu
|
||||
image: diffusers/diffusers-pytorch-cpu
|
||||
report: torch_cpu
|
||||
|
||||
name: ${{ matrix.config.name }}
|
||||
|
||||
runs-on: ${{ matrix.config.runner }}
|
||||
|
||||
container:
|
||||
image: ${{ matrix.config.image }}
|
||||
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
|
||||
|
||||
defaults:
|
||||
run:
|
||||
shell: bash
|
||||
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
uses: actions/checkout@v3
|
||||
with:
|
||||
fetch-depth: 2
|
||||
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
apt-get update && apt-get install libsndfile1-dev -y
|
||||
python -m pip install -e .[quality,test]
|
||||
python -m pip install -U git+https://github.com/huggingface/transformers
|
||||
python -m pip install git+https://github.com/huggingface/accelerate
|
||||
|
||||
- name: Environment
|
||||
run: |
|
||||
python utils/print_env.py
|
||||
|
||||
- name: Run fast PyTorch CPU tests
|
||||
if: ${{ matrix.config.framework == 'pytorch' }}
|
||||
run: |
|
||||
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "not Flax and not Onnx" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Run fast Flax TPU tests
|
||||
if: ${{ matrix.config.framework == 'flax' }}
|
||||
run: |
|
||||
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "Flax" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Run fast ONNXRuntime CPU tests
|
||||
if: ${{ matrix.config.framework == 'onnxruntime' }}
|
||||
run: |
|
||||
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "Onnx" \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
tests/
|
||||
|
||||
- name: Run example PyTorch CPU tests
|
||||
if: ${{ matrix.config.framework == 'pytorch_examples' }}
|
||||
run: |
|
||||
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
|
||||
--make-reports=tests_${{ matrix.config.report }} \
|
||||
examples/test_examples.py
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
|
||||
|
||||
- name: Test suite reports artifacts
|
||||
if: ${{ always() }}
|
||||
uses: actions/upload-artifact@v2
|
||||
with:
|
||||
name: pr_${{ matrix.config.report }}_test_reports
|
||||
path: reports
|
||||
|
||||
run_fast_tests_apple_m1:
|
||||
name: Fast PyTorch MPS tests on MacOS
|
||||
runs-on: [ self-hosted, apple-m1 ]
|
||||
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
uses: actions/checkout@v3
|
||||
with:
|
||||
fetch-depth: 2
|
||||
|
||||
- name: Clean checkout
|
||||
shell: arch -arch arm64 bash {0}
|
||||
run: |
|
||||
git clean -fxd
|
||||
|
||||
- name: Setup miniconda
|
||||
uses: ./.github/actions/setup-miniconda
|
||||
with:
|
||||
python-version: 3.9
|
||||
|
||||
- name: Install dependencies
|
||||
shell: arch -arch arm64 bash {0}
|
||||
run: |
|
||||
${CONDA_RUN} python -m pip install --upgrade pip
|
||||
${CONDA_RUN} python -m pip install -e .[quality,test]
|
||||
${CONDA_RUN} python -m pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
|
||||
${CONDA_RUN} python -m pip install git+https://github.com/huggingface/accelerate
|
||||
${CONDA_RUN} python -m pip install -U git+https://github.com/huggingface/transformers
|
||||
|
||||
- name: Environment
|
||||
shell: arch -arch arm64 bash {0}
|
||||
run: |
|
||||
${CONDA_RUN} python utils/print_env.py
|
||||
|
||||
- name: Run fast PyTorch tests on M1 (MPS)
|
||||
shell: arch -arch arm64 bash {0}
|
||||
env:
|
||||
HF_HOME: /System/Volumes/Data/mnt/cache
|
||||
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
|
||||
run: |
|
||||
${CONDA_RUN} python -m pytest -n 0 -s -v --make-reports=tests_torch_mps tests/
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
run: cat reports/tests_torch_mps_failures_short.txt
|
||||
|
||||
- name: Test suite reports artifacts
|
||||
if: ${{ always() }}
|
||||
uses: actions/upload-artifact@v2
|
||||
with:
|
||||
name: pr_torch_mps_test_reports
|
||||
path: reports
|
||||
10
.gitignore
vendored
10
.gitignore
vendored
@@ -163,12 +163,4 @@ tags
|
||||
*.lock
|
||||
|
||||
# DS_Store (MacOS)
|
||||
.DS_Store
|
||||
# RL pipelines may produce mp4 outputs
|
||||
*.mp4
|
||||
|
||||
# dependencies
|
||||
/transformers
|
||||
|
||||
# ruff
|
||||
.ruff_cache
|
||||
.DS_Store
|
||||
40
CITATION.cff
40
CITATION.cff
@@ -1,40 +0,0 @@
|
||||
cff-version: 1.2.0
|
||||
title: 'Diffusers: State-of-the-art diffusion models'
|
||||
message: >-
|
||||
If you use this software, please cite it using the
|
||||
metadata from this file.
|
||||
type: software
|
||||
authors:
|
||||
- given-names: Patrick
|
||||
family-names: von Platen
|
||||
- given-names: Suraj
|
||||
family-names: Patil
|
||||
- given-names: Anton
|
||||
family-names: Lozhkov
|
||||
- given-names: Pedro
|
||||
family-names: Cuenca
|
||||
- given-names: Nathan
|
||||
family-names: Lambert
|
||||
- given-names: Kashif
|
||||
family-names: Rasul
|
||||
- given-names: Mishig
|
||||
family-names: Davaadorj
|
||||
- given-names: Thomas
|
||||
family-names: Wolf
|
||||
repository-code: 'https://github.com/huggingface/diffusers'
|
||||
abstract: >-
|
||||
Diffusers provides pretrained diffusion models across
|
||||
multiple modalities, such as vision and audio, and serves
|
||||
as a modular toolbox for inference and training of
|
||||
diffusion models.
|
||||
keywords:
|
||||
- deep-learning
|
||||
- pytorch
|
||||
- image-generation
|
||||
- diffusion
|
||||
- text2image
|
||||
- image2image
|
||||
- score-based-generative-modeling
|
||||
- stable-diffusion
|
||||
license: Apache-2.0
|
||||
version: 0.12.1
|
||||
@@ -177,7 +177,7 @@ Follow these steps to start contributing ([supported Python versions](https://gi
|
||||
$ make style
|
||||
```
|
||||
|
||||
🧨 Diffusers also uses `ruff` and a few custom scripts to check for coding mistakes. Quality
|
||||
🧨 Diffusers also uses `flake8` and a few custom scripts to check for coding mistakes. Quality
|
||||
control runs in CI, however you can also run the same checks with:
|
||||
|
||||
```bash
|
||||
|
||||
16
Makefile
16
Makefile
@@ -9,8 +9,9 @@ modified_only_fixup:
|
||||
$(eval modified_py_files := $(shell python utils/get_modified_files.py $(check_dirs)))
|
||||
@if test -n "$(modified_py_files)"; then \
|
||||
echo "Checking/fixing $(modified_py_files)"; \
|
||||
black $(modified_py_files); \
|
||||
ruff $(modified_py_files); \
|
||||
black --preview $(modified_py_files); \
|
||||
isort $(modified_py_files); \
|
||||
flake8 $(modified_py_files); \
|
||||
else \
|
||||
echo "No library .py files were modified"; \
|
||||
fi
|
||||
@@ -40,23 +41,22 @@ repo-consistency:
|
||||
# this target runs checks on all files
|
||||
|
||||
quality:
|
||||
black --check $(check_dirs)
|
||||
ruff $(check_dirs)
|
||||
black --check --preview $(check_dirs)
|
||||
isort --check-only $(check_dirs)
|
||||
flake8 $(check_dirs)
|
||||
doc-builder style src/diffusers docs/source --max_len 119 --check_only --path_to_docs docs/source
|
||||
python utils/check_doc_toc.py
|
||||
|
||||
# Format source code automatically and check is there are any problems left that need manual fixing
|
||||
|
||||
extra_style_checks:
|
||||
python utils/custom_init_isort.py
|
||||
doc-builder style src/diffusers docs/source --max_len 119 --path_to_docs docs/source
|
||||
python utils/check_doc_toc.py --fix_and_overwrite
|
||||
|
||||
# this target runs checks on all files and potentially modifies some of them
|
||||
|
||||
style:
|
||||
black $(check_dirs)
|
||||
ruff $(check_dirs) --fix
|
||||
black --preview $(check_dirs)
|
||||
isort $(check_dirs)
|
||||
${MAKE} autogenerate_code
|
||||
${MAKE} extra_style_checks
|
||||
|
||||
|
||||
200
README.md
200
README.md
@@ -1,6 +1,6 @@
|
||||
<p align="center">
|
||||
<br>
|
||||
<img src="./docs/source/en/imgs/diffusers_library.jpg" width="400"/>
|
||||
<img src="https://github.com/huggingface/diffusers/raw/main/docs/source/imgs/diffusers_library.jpg" width="400"/>
|
||||
<br>
|
||||
<p>
|
||||
<p align="center">
|
||||
@@ -27,28 +27,18 @@ More precisely, 🤗 Diffusers offers:
|
||||
|
||||
## Installation
|
||||
|
||||
### For PyTorch
|
||||
|
||||
**With `pip`** (official package)
|
||||
**With `pip`**
|
||||
|
||||
```bash
|
||||
pip install --upgrade diffusers[torch]
|
||||
pip install --upgrade diffusers
|
||||
```
|
||||
|
||||
**With `conda`** (maintained by the community)
|
||||
**With `conda`**
|
||||
|
||||
```sh
|
||||
conda install -c conda-forge diffusers
|
||||
```
|
||||
|
||||
### For Flax
|
||||
|
||||
**With `pip`**
|
||||
|
||||
```bash
|
||||
pip install --upgrade diffusers[flax]
|
||||
```
|
||||
|
||||
**Apple Silicon (M1/M2) support**
|
||||
|
||||
Please, refer to [the documentation](https://huggingface.co/docs/diffusers/optimization/mps).
|
||||
@@ -79,13 +69,19 @@ In order to get started, we recommend taking a look at two notebooks:
|
||||
Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [LAION](https://laion.ai/) and [RunwayML](https://runwayml.com/). It's trained on 512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/) database. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and runs on a GPU with at least 4GB VRAM.
|
||||
See the [model card](https://huggingface.co/CompVis/stable-diffusion) for more information.
|
||||
|
||||
You need to accept the model license before downloading or using the Stable Diffusion weights. Please, visit the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5), read the license carefully and tick the checkbox if you agree. You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section](https://huggingface.co/docs/hub/security-tokens) of the documentation.
|
||||
|
||||
|
||||
### Text-to-Image generation with Stable Diffusion
|
||||
|
||||
First let's install
|
||||
|
||||
```bash
|
||||
pip install --upgrade diffusers transformers accelerate
|
||||
pip install --upgrade diffusers transformers scipy
|
||||
```
|
||||
|
||||
Run this command to log in with your HF Hub token if you haven't before (you can skip this step if you prefer to run the model locally, follow [this](#running-the-model-locally) instead)
|
||||
```bash
|
||||
huggingface-cli login
|
||||
```
|
||||
|
||||
We recommend using the model in [half-precision (`fp16`)](https://pytorch.org/blog/accelerating-training-on-nvidia-gpus-with-pytorch-automatic-mixed-precision/) as it gives almost always the same results as full
|
||||
@@ -95,7 +91,7 @@ precision while being roughly twice as fast and requiring half the amount of GPU
|
||||
import torch
|
||||
from diffusers import StableDiffusionPipeline
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
|
||||
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, revision="fp16")
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "a photo of an astronaut riding a horse on mars"
|
||||
@@ -103,16 +99,17 @@ image = pipe(prompt).images[0]
|
||||
```
|
||||
|
||||
#### Running the model locally
|
||||
|
||||
You can also simply download the model folder and pass the path to the local folder to the `StableDiffusionPipeline`.
|
||||
If you don't want to login to Hugging Face, you can also simply download the model folder
|
||||
(after having [accepted the license](https://huggingface.co/runwayml/stable-diffusion-v1-5)) and pass
|
||||
the path to the local folder to the `StableDiffusionPipeline`.
|
||||
|
||||
```
|
||||
git lfs install
|
||||
git clone https://huggingface.co/runwayml/stable-diffusion-v1-5
|
||||
```
|
||||
|
||||
Assuming the folder is stored locally under `./stable-diffusion-v1-5`, you can run stable diffusion
|
||||
as follows:
|
||||
Assuming the folder is stored locally under `./stable-diffusion-v1-5`, you can also run stable diffusion
|
||||
without requiring an authentication token:
|
||||
|
||||
```python
|
||||
pipe = StableDiffusionPipeline.from_pretrained("./stable-diffusion-v1-5")
|
||||
@@ -127,7 +124,11 @@ to using `fp16`.
|
||||
The following snippet should result in less than 4GB VRAM.
|
||||
|
||||
```python
|
||||
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
|
||||
pipe = StableDiffusionPipeline.from_pretrained(
|
||||
"runwayml/stable-diffusion-v1-5",
|
||||
revision="fp16",
|
||||
torch_dtype=torch.float16,
|
||||
)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "a photo of an astronaut riding a horse on mars"
|
||||
@@ -141,7 +142,19 @@ it before the pipeline and pass it to `from_pretrained`.
|
||||
```python
|
||||
from diffusers import LMSDiscreteScheduler
|
||||
|
||||
pipe.scheduler = LMSDiscreteScheduler.from_config(pipe.scheduler.config)
|
||||
lms = LMSDiscreteScheduler(
|
||||
beta_start=0.00085,
|
||||
beta_end=0.012,
|
||||
beta_schedule="scaled_linear"
|
||||
)
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained(
|
||||
"runwayml/stable-diffusion-v1-5",
|
||||
revision="fp16",
|
||||
torch_dtype=torch.float16,
|
||||
scheduler=lms,
|
||||
)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "a photo of an astronaut riding a horse on mars"
|
||||
image = pipe(prompt).images[0]
|
||||
@@ -153,6 +166,7 @@ If you want to run Stable Diffusion on CPU or you want to have maximum precision
|
||||
please run the model in the default *full-precision* setting:
|
||||
|
||||
```python
|
||||
# make sure you're logged in with `huggingface-cli login`
|
||||
from diffusers import StableDiffusionPipeline
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
|
||||
@@ -235,102 +249,6 @@ images = pipeline(prompt_ids, params, prng_seed, num_inference_steps, jit=True).
|
||||
images = pipeline.numpy_to_pil(np.asarray(images.reshape((num_samples,) + images.shape[-3:])))
|
||||
```
|
||||
|
||||
Diffusers also has a Image-to-Image generation pipeline with Flax/Jax
|
||||
```python
|
||||
import jax
|
||||
import numpy as np
|
||||
import jax.numpy as jnp
|
||||
from flax.jax_utils import replicate
|
||||
from flax.training.common_utils import shard
|
||||
import requests
|
||||
from io import BytesIO
|
||||
from PIL import Image
|
||||
from diffusers import FlaxStableDiffusionImg2ImgPipeline
|
||||
|
||||
def create_key(seed=0):
|
||||
return jax.random.PRNGKey(seed)
|
||||
rng = create_key(0)
|
||||
|
||||
url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
|
||||
response = requests.get(url)
|
||||
init_img = Image.open(BytesIO(response.content)).convert("RGB")
|
||||
init_img = init_img.resize((768, 512))
|
||||
|
||||
prompts = "A fantasy landscape, trending on artstation"
|
||||
|
||||
pipeline, params = FlaxStableDiffusionImg2ImgPipeline.from_pretrained(
|
||||
"CompVis/stable-diffusion-v1-4", revision="flax",
|
||||
dtype=jnp.bfloat16,
|
||||
)
|
||||
|
||||
num_samples = jax.device_count()
|
||||
rng = jax.random.split(rng, jax.device_count())
|
||||
prompt_ids, processed_image = pipeline.prepare_inputs(prompt=[prompts]*num_samples, image = [init_img]*num_samples)
|
||||
p_params = replicate(params)
|
||||
prompt_ids = shard(prompt_ids)
|
||||
processed_image = shard(processed_image)
|
||||
|
||||
output = pipeline(
|
||||
prompt_ids=prompt_ids,
|
||||
image=processed_image,
|
||||
params=p_params,
|
||||
prng_seed=rng,
|
||||
strength=0.75,
|
||||
num_inference_steps=50,
|
||||
jit=True,
|
||||
height=512,
|
||||
width=768).images
|
||||
|
||||
output_images = pipeline.numpy_to_pil(np.asarray(output.reshape((num_samples,) + output.shape[-3:])))
|
||||
```
|
||||
|
||||
Diffusers also has a Text-guided inpainting pipeline with Flax/Jax
|
||||
|
||||
```python
|
||||
import jax
|
||||
import numpy as np
|
||||
from flax.jax_utils import replicate
|
||||
from flax.training.common_utils import shard
|
||||
import PIL
|
||||
import requests
|
||||
from io import BytesIO
|
||||
|
||||
|
||||
from diffusers import FlaxStableDiffusionInpaintPipeline
|
||||
|
||||
def download_image(url):
|
||||
response = requests.get(url)
|
||||
return PIL.Image.open(BytesIO(response.content)).convert("RGB")
|
||||
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
|
||||
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
|
||||
|
||||
init_image = download_image(img_url).resize((512, 512))
|
||||
mask_image = download_image(mask_url).resize((512, 512))
|
||||
|
||||
pipeline, params = FlaxStableDiffusionInpaintPipeline.from_pretrained("xvjiarui/stable-diffusion-2-inpainting")
|
||||
|
||||
prompt = "Face of a yellow cat, high resolution, sitting on a park bench"
|
||||
prng_seed = jax.random.PRNGKey(0)
|
||||
num_inference_steps = 50
|
||||
|
||||
num_samples = jax.device_count()
|
||||
prompt = num_samples * [prompt]
|
||||
init_image = num_samples * [init_image]
|
||||
mask_image = num_samples * [mask_image]
|
||||
prompt_ids, processed_masked_images, processed_masks = pipeline.prepare_inputs(prompt, init_image, mask_image)
|
||||
|
||||
|
||||
# shard inputs and rng
|
||||
params = replicate(params)
|
||||
prng_seed = jax.random.split(prng_seed, jax.device_count())
|
||||
prompt_ids = shard(prompt_ids)
|
||||
processed_masked_images = shard(processed_masked_images)
|
||||
processed_masks = shard(processed_masks)
|
||||
|
||||
images = pipeline(prompt_ids, processed_masks, processed_masked_images, params, prng_seed, num_inference_steps, jit=True).images
|
||||
images = pipeline.numpy_to_pil(np.asarray(images.reshape((num_samples,) + images.shape[-3:])))
|
||||
```
|
||||
|
||||
### Image-to-Image text-guided generation with Stable Diffusion
|
||||
|
||||
The `StableDiffusionImg2ImgPipeline` lets you pass a text prompt and an initial image to condition the generation of new images.
|
||||
@@ -346,8 +264,11 @@ from diffusers import StableDiffusionImg2ImgPipeline
|
||||
# load the pipeline
|
||||
device = "cuda"
|
||||
model_id_or_path = "runwayml/stable-diffusion-v1-5"
|
||||
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(model_id_or_path, torch_dtype=torch.float16)
|
||||
|
||||
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
|
||||
model_id_or_path,
|
||||
revision="fp16",
|
||||
torch_dtype=torch.float16,
|
||||
)
|
||||
# or download via git clone https://huggingface.co/runwayml/stable-diffusion-v1-5
|
||||
# and pass `model_id_or_path="./stable-diffusion-v1-5"`.
|
||||
pipe = pipe.to(device)
|
||||
@@ -361,7 +282,7 @@ init_image = init_image.resize((768, 512))
|
||||
|
||||
prompt = "A fantasy landscape, trending on artstation"
|
||||
|
||||
images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
|
||||
images[0].save("fantasy_landscape.png")
|
||||
```
|
||||
@@ -369,7 +290,10 @@ You can also run this example on colab [, read the license carefully and tick the checkbox if you agree. Note that this is an additional license, you need to accept it even if you accepted the text-to-image Stable Diffusion license in the past. You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section](https://huggingface.co/docs/hub/security-tokens) of the documentation.
|
||||
|
||||
|
||||
```python
|
||||
import PIL
|
||||
@@ -389,7 +313,11 @@ mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data
|
||||
init_image = download_image(img_url).resize((512, 512))
|
||||
mask_image = download_image(mask_url).resize((512, 512))
|
||||
|
||||
pipe = StableDiffusionInpaintPipeline.from_pretrained("runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16)
|
||||
pipe = StableDiffusionInpaintPipeline.from_pretrained(
|
||||
"runwayml/stable-diffusion-inpainting",
|
||||
revision="fp16",
|
||||
torch_dtype=torch.float16,
|
||||
)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "Face of a yellow cat, high resolution, sitting on a park bench"
|
||||
@@ -398,8 +326,11 @@ image = pipe(prompt=prompt, image=init_image, mask_image=mask_image).images[0]
|
||||
|
||||
### Tweak prompts reusing seeds and latents
|
||||
|
||||
You can generate your own latents to reproduce results, or tweak your prompt on a specific result you liked.
|
||||
Please have a look at [Reusing seeds for deterministic generation](https://huggingface.co/docs/diffusers/main/en/using-diffusers/reusing_seeds).
|
||||
You can generate your own latents to reproduce results, or tweak your prompt on a specific result you liked. [This notebook](https://github.com/pcuenca/diffusers-examples/blob/main/notebooks/stable-diffusion-seeds.ipynb) shows how to do it step by step. You can also run it in Google Colab [](https://colab.research.google.com/github/pcuenca/diffusers-examples/blob/main/notebooks/stable-diffusion-seeds.ipynb).
|
||||
|
||||
|
||||
For more details, check out [the Stable Diffusion notebook](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb) [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb)
|
||||
and have a look into the [release notes](https://github.com/huggingface/diffusers/releases/tag/v0.2.0).
|
||||
|
||||
## Fine-Tuning Stable Diffusion
|
||||
|
||||
@@ -409,15 +340,14 @@ Textual Inversion is a technique for capturing novel concepts from a small numbe
|
||||
|
||||
- Textual Inversion. Capture novel concepts from a small set of sample images, and associate them with new "words" in the embedding space of the text encoder. Please, refer to [our training examples](https://github.com/huggingface/diffusers/tree/main/examples/textual_inversion) or [documentation](https://huggingface.co/docs/diffusers/training/text_inversion) to try for yourself.
|
||||
|
||||
- Dreambooth. Another technique to capture new concepts in Stable Diffusion. This method fine-tunes the UNet (and, optionally, also the text encoder) of the pipeline to achieve impressive results. Please, refer to [our training example](https://github.com/huggingface/diffusers/tree/main/examples/dreambooth) and [training report](https://huggingface.co/blog/dreambooth) for additional details and training recommendations.
|
||||
- Dreambooth. Another technique to capture new concepts in Stable Diffusion. This method fine-tunes the UNet (and, optionally, also the text encoder) of the pipeline to achieve impressive results. Please, refer to [our training examples](https://github.com/huggingface/diffusers/tree/main/examples/dreambooth) and [training report](https://wandb.ai/psuraj/dreambooth/reports/Dreambooth-Training-Analysis--VmlldzoyNzk0NDc3) for additional details and training recommendations.
|
||||
|
||||
- Full Stable Diffusion fine-tuning. If you have a more sizable dataset with a specific look or style, you can fine-tune Stable Diffusion so that it outputs images following those examples. This was the approach taken to create [a Pokémon Stable Diffusion model](https://huggingface.co/justinpinkney/pokemon-stable-diffusion) (by Justing Pinkney / Lambda Labs), [a Japanese specific version of Stable Diffusion](https://huggingface.co/spaces/rinna/japanese-stable-diffusion) (by [Rinna Co.](https://github.com/rinnakk/japanese-stable-diffusion/) and others. You can start at [our text-to-image fine-tuning example](https://github.com/huggingface/diffusers/tree/main/examples/text_to_image) and go from there.
|
||||
|
||||
|
||||
## Stable Diffusion Community Pipelines
|
||||
|
||||
The release of Stable Diffusion as an open source model has fostered a lot of interesting ideas and experimentation.
|
||||
Our [Community Examples folder](https://github.com/huggingface/diffusers/tree/main/examples/community) contains many ideas worth exploring, like interpolating to create animated videos, using CLIP Guidance for additional prompt fidelity, term weighting, and much more! [Take a look](https://huggingface.co/docs/diffusers/using-diffusers/custom_pipeline_overview) and [contribute your own](https://huggingface.co/docs/diffusers/using-diffusers/contribute_pipeline).
|
||||
The release of Stable Diffusion as an open source model has fostered a lot of interesting ideas and experimentation. Our [Community Examples folder](https://github.com/huggingface/diffusers/tree/main/examples/community) contains many ideas worth exploring, like interpolating to create animated videos, using CLIP Guidance for additional prompt fidelity, term weighting, and much more! Take a look and [contribute your own](https://huggingface.co/docs/diffusers/using-diffusers/custom_pipelines).
|
||||
|
||||
## Other Examples
|
||||
|
||||
@@ -428,7 +358,7 @@ There are many ways to try running Diffusers! Here we outline code-focused tools
|
||||
If you want to run the code yourself 💻, you can try out:
|
||||
- [Text-to-Image Latent Diffusion](https://huggingface.co/CompVis/ldm-text2im-large-256)
|
||||
```python
|
||||
# !pip install diffusers["torch"] transformers
|
||||
# !pip install diffusers transformers
|
||||
from diffusers import DiffusionPipeline
|
||||
|
||||
device = "cuda"
|
||||
@@ -447,7 +377,7 @@ image.save("squirrel.png")
|
||||
```
|
||||
- [Unconditional Diffusion with discrete scheduler](https://huggingface.co/google/ddpm-celebahq-256)
|
||||
```python
|
||||
# !pip install diffusers["torch"]
|
||||
# !pip install diffusers
|
||||
from diffusers import DDPMPipeline, DDIMPipeline, PNDMPipeline
|
||||
|
||||
model_id = "google/ddpm-celebahq-256"
|
||||
@@ -466,14 +396,10 @@ image.save("ddpm_generated_image.png")
|
||||
- [Unconditional Latent Diffusion](https://huggingface.co/CompVis/ldm-celebahq-256)
|
||||
- [Unconditional Diffusion with continuous scheduler](https://huggingface.co/google/ncsnpp-ffhq-1024)
|
||||
|
||||
**Other Image Notebooks**:
|
||||
**Other Notebooks**:
|
||||
* [image-to-image generation with Stable Diffusion](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb) ,
|
||||
* [tweak images via repeated Stable Diffusion seeds](https://colab.research.google.com/github/pcuenca/diffusers-examples/blob/main/notebooks/stable-diffusion-seeds.ipynb) ,
|
||||
|
||||
**Diffusers for Other Modalities**:
|
||||
* [Molecule conformation generation](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/geodiff_molecule_conformation.ipynb) ,
|
||||
* [Model-based reinforcement learning](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/reinforcement_learning_with_diffusers.ipynb) ,
|
||||
|
||||
### Web Demos
|
||||
If you just want to play around with some web demos, you can try out the following 🚀 Spaces:
|
||||
| Model | Hugging Face Spaces |
|
||||
@@ -496,7 +422,7 @@ If you just want to play around with some web demos, you can try out the followi
|
||||
<p>
|
||||
|
||||
**Schedulers**: Algorithm class for both **inference** and **training**.
|
||||
The class provides functionality to compute previous image according to alpha, beta schedule as well as predict noise for training. Also known as **Samplers**.
|
||||
The class provides functionality to compute previous image according to alpha, beta schedule as well as predict noise for training.
|
||||
*Examples*: [DDPM](https://arxiv.org/abs/2006.11239), [DDIM](https://arxiv.org/abs/2010.02502), [PNDM](https://arxiv.org/abs/2202.09778), [DEIS](https://arxiv.org/abs/2204.13902)
|
||||
|
||||
<p align="center">
|
||||
|
||||
@@ -1,44 +0,0 @@
|
||||
FROM ubuntu:20.04
|
||||
LABEL maintainer="Hugging Face"
|
||||
LABEL repository="diffusers"
|
||||
|
||||
ENV DEBIAN_FRONTEND=noninteractive
|
||||
|
||||
RUN apt update && \
|
||||
apt install -y bash \
|
||||
build-essential \
|
||||
git \
|
||||
git-lfs \
|
||||
curl \
|
||||
ca-certificates \
|
||||
libsndfile1-dev \
|
||||
python3.8 \
|
||||
python3-pip \
|
||||
python3.8-venv && \
|
||||
rm -rf /var/lib/apt/lists
|
||||
|
||||
# make sure to use venv
|
||||
RUN python3 -m venv /opt/venv
|
||||
ENV PATH="/opt/venv/bin:$PATH"
|
||||
|
||||
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
|
||||
# follow the instructions here: https://cloud.google.com/tpu/docs/run-in-container#train_a_jax_model_in_a_docker_container
|
||||
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
|
||||
python3 -m pip install --upgrade --no-cache-dir \
|
||||
clu \
|
||||
"jax[cpu]>=0.2.16,!=0.3.2" \
|
||||
"flax>=0.4.1" \
|
||||
"jaxlib>=0.1.65" && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
accelerate \
|
||||
datasets \
|
||||
hf-doc-builder \
|
||||
huggingface-hub \
|
||||
Jinja2 \
|
||||
librosa \
|
||||
numpy \
|
||||
scipy \
|
||||
tensorboard \
|
||||
transformers
|
||||
|
||||
CMD ["/bin/bash"]
|
||||
@@ -1,46 +0,0 @@
|
||||
FROM ubuntu:20.04
|
||||
LABEL maintainer="Hugging Face"
|
||||
LABEL repository="diffusers"
|
||||
|
||||
ENV DEBIAN_FRONTEND=noninteractive
|
||||
|
||||
RUN apt update && \
|
||||
apt install -y bash \
|
||||
build-essential \
|
||||
git \
|
||||
git-lfs \
|
||||
curl \
|
||||
ca-certificates \
|
||||
libsndfile1-dev \
|
||||
python3.8 \
|
||||
python3-pip \
|
||||
python3.8-venv && \
|
||||
rm -rf /var/lib/apt/lists
|
||||
|
||||
# make sure to use venv
|
||||
RUN python3 -m venv /opt/venv
|
||||
ENV PATH="/opt/venv/bin:$PATH"
|
||||
|
||||
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
|
||||
# follow the instructions here: https://cloud.google.com/tpu/docs/run-in-container#train_a_jax_model_in_a_docker_container
|
||||
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
"jax[tpu]>=0.2.16,!=0.3.2" \
|
||||
-f https://storage.googleapis.com/jax-releases/libtpu_releases.html && \
|
||||
python3 -m pip install --upgrade --no-cache-dir \
|
||||
clu \
|
||||
"flax>=0.4.1" \
|
||||
"jaxlib>=0.1.65" && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
accelerate \
|
||||
datasets \
|
||||
hf-doc-builder \
|
||||
huggingface-hub \
|
||||
Jinja2 \
|
||||
librosa \
|
||||
numpy \
|
||||
scipy \
|
||||
tensorboard \
|
||||
transformers
|
||||
|
||||
CMD ["/bin/bash"]
|
||||
@@ -1,44 +0,0 @@
|
||||
FROM ubuntu:20.04
|
||||
LABEL maintainer="Hugging Face"
|
||||
LABEL repository="diffusers"
|
||||
|
||||
ENV DEBIAN_FRONTEND=noninteractive
|
||||
|
||||
RUN apt update && \
|
||||
apt install -y bash \
|
||||
build-essential \
|
||||
git \
|
||||
git-lfs \
|
||||
curl \
|
||||
ca-certificates \
|
||||
libsndfile1-dev \
|
||||
python3.8 \
|
||||
python3-pip \
|
||||
python3.8-venv && \
|
||||
rm -rf /var/lib/apt/lists
|
||||
|
||||
# make sure to use venv
|
||||
RUN python3 -m venv /opt/venv
|
||||
ENV PATH="/opt/venv/bin:$PATH"
|
||||
|
||||
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
|
||||
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
torch \
|
||||
torchvision \
|
||||
torchaudio \
|
||||
onnxruntime \
|
||||
--extra-index-url https://download.pytorch.org/whl/cpu && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
accelerate \
|
||||
datasets \
|
||||
hf-doc-builder \
|
||||
huggingface-hub \
|
||||
Jinja2 \
|
||||
librosa \
|
||||
numpy \
|
||||
scipy \
|
||||
tensorboard \
|
||||
transformers
|
||||
|
||||
CMD ["/bin/bash"]
|
||||
@@ -1,44 +0,0 @@
|
||||
FROM nvidia/cuda:11.6.2-cudnn8-devel-ubuntu20.04
|
||||
LABEL maintainer="Hugging Face"
|
||||
LABEL repository="diffusers"
|
||||
|
||||
ENV DEBIAN_FRONTEND=noninteractive
|
||||
|
||||
RUN apt update && \
|
||||
apt install -y bash \
|
||||
build-essential \
|
||||
git \
|
||||
git-lfs \
|
||||
curl \
|
||||
ca-certificates \
|
||||
libsndfile1-dev \
|
||||
python3.8 \
|
||||
python3-pip \
|
||||
python3.8-venv && \
|
||||
rm -rf /var/lib/apt/lists
|
||||
|
||||
# make sure to use venv
|
||||
RUN python3 -m venv /opt/venv
|
||||
ENV PATH="/opt/venv/bin:$PATH"
|
||||
|
||||
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
|
||||
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
torch \
|
||||
torchvision \
|
||||
torchaudio \
|
||||
"onnxruntime-gpu>=1.13.1" \
|
||||
--extra-index-url https://download.pytorch.org/whl/cu117 && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
accelerate \
|
||||
datasets \
|
||||
hf-doc-builder \
|
||||
huggingface-hub \
|
||||
Jinja2 \
|
||||
librosa \
|
||||
numpy \
|
||||
scipy \
|
||||
tensorboard \
|
||||
transformers
|
||||
|
||||
CMD ["/bin/bash"]
|
||||
@@ -1,43 +0,0 @@
|
||||
FROM ubuntu:20.04
|
||||
LABEL maintainer="Hugging Face"
|
||||
LABEL repository="diffusers"
|
||||
|
||||
ENV DEBIAN_FRONTEND=noninteractive
|
||||
|
||||
RUN apt update && \
|
||||
apt install -y bash \
|
||||
build-essential \
|
||||
git \
|
||||
git-lfs \
|
||||
curl \
|
||||
ca-certificates \
|
||||
libsndfile1-dev \
|
||||
python3.8 \
|
||||
python3-pip \
|
||||
python3.8-venv && \
|
||||
rm -rf /var/lib/apt/lists
|
||||
|
||||
# make sure to use venv
|
||||
RUN python3 -m venv /opt/venv
|
||||
ENV PATH="/opt/venv/bin:$PATH"
|
||||
|
||||
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
|
||||
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
torch \
|
||||
torchvision \
|
||||
torchaudio \
|
||||
--extra-index-url https://download.pytorch.org/whl/cpu && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
accelerate \
|
||||
datasets \
|
||||
hf-doc-builder \
|
||||
huggingface-hub \
|
||||
Jinja2 \
|
||||
librosa \
|
||||
numpy \
|
||||
scipy \
|
||||
tensorboard \
|
||||
transformers
|
||||
|
||||
CMD ["/bin/bash"]
|
||||
@@ -1,43 +0,0 @@
|
||||
FROM nvidia/cuda:11.7.1-cudnn8-runtime-ubuntu20.04
|
||||
LABEL maintainer="Hugging Face"
|
||||
LABEL repository="diffusers"
|
||||
|
||||
ENV DEBIAN_FRONTEND=noninteractive
|
||||
|
||||
RUN apt update && \
|
||||
apt install -y bash \
|
||||
build-essential \
|
||||
git \
|
||||
git-lfs \
|
||||
curl \
|
||||
ca-certificates \
|
||||
libsndfile1-dev \
|
||||
python3.8 \
|
||||
python3-pip \
|
||||
python3.8-venv && \
|
||||
rm -rf /var/lib/apt/lists
|
||||
|
||||
# make sure to use venv
|
||||
RUN python3 -m venv /opt/venv
|
||||
ENV PATH="/opt/venv/bin:$PATH"
|
||||
|
||||
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
|
||||
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
torch \
|
||||
torchvision \
|
||||
torchaudio \
|
||||
--extra-index-url https://download.pytorch.org/whl/cu117 && \
|
||||
python3 -m pip install --no-cache-dir \
|
||||
accelerate \
|
||||
datasets \
|
||||
hf-doc-builder \
|
||||
huggingface-hub \
|
||||
Jinja2 \
|
||||
librosa \
|
||||
numpy \
|
||||
scipy \
|
||||
tensorboard \
|
||||
transformers
|
||||
|
||||
CMD ["/bin/bash"]
|
||||
271
docs/README.md
271
docs/README.md
@@ -1,271 +0,0 @@
|
||||
<!---
|
||||
Copyright 2022- The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License");
|
||||
you may not use this file except in compliance with the License.
|
||||
You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software
|
||||
distributed under the License is distributed on an "AS IS" BASIS,
|
||||
WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
See the License for the specific language governing permissions and
|
||||
limitations under the License.
|
||||
-->
|
||||
|
||||
# Generating the documentation
|
||||
|
||||
To generate the documentation, you first have to build it. Several packages are necessary to build the doc,
|
||||
you can install them with the following command, at the root of the code repository:
|
||||
|
||||
```bash
|
||||
pip install -e ".[docs]"
|
||||
```
|
||||
|
||||
Then you need to install our open source documentation builder tool:
|
||||
|
||||
```bash
|
||||
pip install git+https://github.com/huggingface/doc-builder
|
||||
```
|
||||
|
||||
---
|
||||
**NOTE**
|
||||
|
||||
You only need to generate the documentation to inspect it locally (if you're planning changes and want to
|
||||
check how they look before committing for instance). You don't have to commit the built documentation.
|
||||
|
||||
---
|
||||
|
||||
## Previewing the documentation
|
||||
|
||||
To preview the docs, first install the `watchdog` module with:
|
||||
|
||||
```bash
|
||||
pip install watchdog
|
||||
```
|
||||
|
||||
Then run the following command:
|
||||
|
||||
```bash
|
||||
doc-builder preview {package_name} {path_to_docs}
|
||||
```
|
||||
|
||||
For example:
|
||||
|
||||
```bash
|
||||
doc-builder preview diffusers docs/source/en
|
||||
```
|
||||
|
||||
The docs will be viewable at [http://localhost:3000](http://localhost:3000). You can also preview the docs once you have opened a PR. You will see a bot add a comment to a link where the documentation with your changes lives.
|
||||
|
||||
---
|
||||
**NOTE**
|
||||
|
||||
The `preview` command only works with existing doc files. When you add a completely new file, you need to update `_toctree.yml` & restart `preview` command (`ctrl-c` to stop it & call `doc-builder preview ...` again).
|
||||
|
||||
---
|
||||
|
||||
## Adding a new element to the navigation bar
|
||||
|
||||
Accepted files are Markdown (.md or .mdx).
|
||||
|
||||
Create a file with its extension and put it in the source directory. You can then link it to the toc-tree by putting
|
||||
the filename without the extension in the [`_toctree.yml`](https://github.com/huggingface/diffusers/blob/main/docs/source/_toctree.yml) file.
|
||||
|
||||
## Renaming section headers and moving sections
|
||||
|
||||
It helps to keep the old links working when renaming the section header and/or moving sections from one document to another. This is because the old links are likely to be used in Issues, Forums, and Social media and it'd make for a much more superior user experience if users reading those months later could still easily navigate to the originally intended information.
|
||||
|
||||
Therefore, we simply keep a little map of moved sections at the end of the document where the original section was. The key is to preserve the original anchor.
|
||||
|
||||
So if you renamed a section from: "Section A" to "Section B", then you can add at the end of the file:
|
||||
|
||||
```
|
||||
Sections that were moved:
|
||||
|
||||
[ <a href="#section-b">Section A</a><a id="section-a"></a> ]
|
||||
```
|
||||
and of course, if you moved it to another file, then:
|
||||
|
||||
```
|
||||
Sections that were moved:
|
||||
|
||||
[ <a href="../new-file#section-b">Section A</a><a id="section-a"></a> ]
|
||||
```
|
||||
|
||||
Use the relative style to link to the new file so that the versioned docs continue to work.
|
||||
|
||||
For an example of a rich moved section set please see the very end of [the transformers Trainer doc](https://github.com/huggingface/transformers/blob/main/docs/source/en/main_classes/trainer.mdx).
|
||||
|
||||
|
||||
## Writing Documentation - Specification
|
||||
|
||||
The `huggingface/diffusers` documentation follows the
|
||||
[Google documentation](https://sphinxcontrib-napoleon.readthedocs.io/en/latest/example_google.html) style for docstrings,
|
||||
although we can write them directly in Markdown.
|
||||
|
||||
### Adding a new tutorial
|
||||
|
||||
Adding a new tutorial or section is done in two steps:
|
||||
|
||||
- Add a new file under `docs/source`. This file can either be ReStructuredText (.rst) or Markdown (.md).
|
||||
- Link that file in `docs/source/_toctree.yml` on the correct toc-tree.
|
||||
|
||||
Make sure to put your new file under the proper section. It's unlikely to go in the first section (*Get Started*), so
|
||||
depending on the intended targets (beginners, more advanced users, or researchers) it should go in sections two, three, or four.
|
||||
|
||||
### Adding a new pipeline/scheduler
|
||||
|
||||
When adding a new pipeline:
|
||||
|
||||
- create a file `xxx.mdx` under `docs/source/api/pipelines` (don't hesitate to copy an existing file as template).
|
||||
- Link that file in (*Diffusers Summary*) section in `docs/source/api/pipelines/overview.mdx`, along with the link to the paper, and a colab notebook (if available).
|
||||
- Write a short overview of the diffusion model:
|
||||
- Overview with paper & authors
|
||||
- Paper abstract
|
||||
- Tips and tricks and how to use it best
|
||||
- Possible an end-to-end example of how to use it
|
||||
- Add all the pipeline classes that should be linked in the diffusion model. These classes should be added using our Markdown syntax. By default as follows:
|
||||
|
||||
```
|
||||
## XXXPipeline
|
||||
|
||||
[[autodoc]] XXXPipeline
|
||||
- all
|
||||
- __call__
|
||||
```
|
||||
|
||||
This will include every public method of the pipeline that is documented, as well as the `__call__` method that is not documented by default. If you just want to add additional methods that are not documented, you can put the list of all methods to add in a list that contains `all`.
|
||||
|
||||
```
|
||||
[[autodoc]] XXXPipeline
|
||||
- all
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
```
|
||||
|
||||
You can follow the same process to create a new scheduler under the `docs/source/api/schedulers` folder
|
||||
|
||||
### Writing source documentation
|
||||
|
||||
Values that should be put in `code` should either be surrounded by backticks: \`like so\`. Note that argument names
|
||||
and objects like True, None, or any strings should usually be put in `code`.
|
||||
|
||||
When mentioning a class, function, or method, it is recommended to use our syntax for internal links so that our tool
|
||||
adds a link to its documentation with this syntax: \[\`XXXClass\`\] or \[\`function\`\]. This requires the class or
|
||||
function to be in the main package.
|
||||
|
||||
If you want to create a link to some internal class or function, you need to
|
||||
provide its path. For instance: \[\`pipelines.ImagePipelineOutput\`\]. This will be converted into a link with
|
||||
`pipelines.ImagePipelineOutput` in the description. To get rid of the path and only keep the name of the object you are
|
||||
linking to in the description, add a ~: \[\`~pipelines.ImagePipelineOutput\`\] will generate a link with `ImagePipelineOutput` in the description.
|
||||
|
||||
The same works for methods so you can either use \[\`XXXClass.method\`\] or \[~\`XXXClass.method\`\].
|
||||
|
||||
#### Defining arguments in a method
|
||||
|
||||
Arguments should be defined with the `Args:` (or `Arguments:` or `Parameters:`) prefix, followed by a line return and
|
||||
an indentation. The argument should be followed by its type, with its shape if it is a tensor, a colon, and its
|
||||
description:
|
||||
|
||||
```
|
||||
Args:
|
||||
n_layers (`int`): The number of layers of the model.
|
||||
```
|
||||
|
||||
If the description is too long to fit in one line, another indentation is necessary before writing the description
|
||||
after the argument.
|
||||
|
||||
Here's an example showcasing everything so far:
|
||||
|
||||
```
|
||||
Args:
|
||||
input_ids (`torch.LongTensor` of shape `(batch_size, sequence_length)`):
|
||||
Indices of input sequence tokens in the vocabulary.
|
||||
|
||||
Indices can be obtained using [`AlbertTokenizer`]. See [`~PreTrainedTokenizer.encode`] and
|
||||
[`~PreTrainedTokenizer.__call__`] for details.
|
||||
|
||||
[What are input IDs?](../glossary#input-ids)
|
||||
```
|
||||
|
||||
For optional arguments or arguments with defaults we follow the following syntax: imagine we have a function with the
|
||||
following signature:
|
||||
|
||||
```
|
||||
def my_function(x: str = None, a: float = 1):
|
||||
```
|
||||
|
||||
then its documentation should look like this:
|
||||
|
||||
```
|
||||
Args:
|
||||
x (`str`, *optional*):
|
||||
This argument controls ...
|
||||
a (`float`, *optional*, defaults to 1):
|
||||
This argument is used to ...
|
||||
```
|
||||
|
||||
Note that we always omit the "defaults to \`None\`" when None is the default for any argument. Also note that even
|
||||
if the first line describing your argument type and its default gets long, you can't break it on several lines. You can
|
||||
however write as many lines as you want in the indented description (see the example above with `input_ids`).
|
||||
|
||||
#### Writing a multi-line code block
|
||||
|
||||
Multi-line code blocks can be useful for displaying examples. They are done between two lines of three backticks as usual in Markdown:
|
||||
|
||||
|
||||
````
|
||||
```
|
||||
# first line of code
|
||||
# second line
|
||||
# etc
|
||||
```
|
||||
````
|
||||
|
||||
#### Writing a return block
|
||||
|
||||
The return block should be introduced with the `Returns:` prefix, followed by a line return and an indentation.
|
||||
The first line should be the type of the return, followed by a line return. No need to indent further for the elements
|
||||
building the return.
|
||||
|
||||
Here's an example of a single value return:
|
||||
|
||||
```
|
||||
Returns:
|
||||
`List[int]`: A list of integers in the range [0, 1] --- 1 for a special token, 0 for a sequence token.
|
||||
```
|
||||
|
||||
Here's an example of a tuple return, comprising several objects:
|
||||
|
||||
```
|
||||
Returns:
|
||||
`tuple(torch.FloatTensor)` comprising various elements depending on the configuration ([`BertConfig`]) and inputs:
|
||||
- ** loss** (*optional*, returned when `masked_lm_labels` is provided) `torch.FloatTensor` of shape `(1,)` --
|
||||
Total loss is the sum of the masked language modeling loss and the next sequence prediction (classification) loss.
|
||||
- **prediction_scores** (`torch.FloatTensor` of shape `(batch_size, sequence_length, config.vocab_size)`) --
|
||||
Prediction scores of the language modeling head (scores for each vocabulary token before SoftMax).
|
||||
```
|
||||
|
||||
#### Adding an image
|
||||
|
||||
Due to the rapidly growing repository, it is important to make sure that no files that would significantly weigh down the repository are added. This includes images, videos, and other non-text files. We prefer to leverage a hf.co hosted `dataset` like
|
||||
the ones hosted on [`hf-internal-testing`](https://huggingface.co/hf-internal-testing) in which to place these files and reference
|
||||
them by URL. We recommend putting them in the following dataset: [huggingface/documentation-images](https://huggingface.co/datasets/huggingface/documentation-images).
|
||||
If an external contribution, feel free to add the images to your PR and ask a Hugging Face member to migrate your images
|
||||
to this dataset.
|
||||
|
||||
## Styling the docstring
|
||||
|
||||
We have an automatic script running with the `make style` command that will make sure that:
|
||||
- the docstrings fully take advantage of the line width
|
||||
- all code examples are formatted using black, like the code of the Transformers library
|
||||
|
||||
This script may have some weird failures if you made a syntax mistake or if you uncover a bug. Therefore, it's
|
||||
recommended to commit your changes before running `make style`, so you can revert the changes done by that script
|
||||
easily.
|
||||
|
||||
@@ -1,57 +0,0 @@
|
||||
### Translating the Diffusers documentation into your language
|
||||
|
||||
As part of our mission to democratize machine learning, we'd love to make the Diffusers library available in many more languages! Follow the steps below if you want to help translate the documentation into your language 🙏.
|
||||
|
||||
**🗞️ Open an issue**
|
||||
|
||||
To get started, navigate to the [Issues](https://github.com/huggingface/diffusers/issues) page of this repo and check if anyone else has opened an issue for your language. If not, open a new issue by selecting the "Translation template" from the "New issue" button.
|
||||
|
||||
Once an issue exists, post a comment to indicate which chapters you'd like to work on, and we'll add your name to the list.
|
||||
|
||||
|
||||
**🍴 Fork the repository**
|
||||
|
||||
First, you'll need to [fork the Diffusers repo](https://docs.github.com/en/get-started/quickstart/fork-a-repo). You can do this by clicking on the **Fork** button on the top-right corner of this repo's page.
|
||||
|
||||
Once you've forked the repo, you'll want to get the files on your local machine for editing. You can do that by cloning the fork with Git as follows:
|
||||
|
||||
```bash
|
||||
git clone https://github.com/YOUR-USERNAME/diffusers.git
|
||||
```
|
||||
|
||||
**📋 Copy-paste the English version with a new language code**
|
||||
|
||||
The documentation files are in one leading directory:
|
||||
|
||||
- [`docs/source`](https://github.com/huggingface/diffusers/tree/main/docs/source): All the documentation materials are organized here by language.
|
||||
|
||||
You'll only need to copy the files in the [`docs/source/en`](https://github.com/huggingface/diffusers/tree/main/docs/source/en) directory, so first navigate to your fork of the repo and run the following:
|
||||
|
||||
```bash
|
||||
cd ~/path/to/diffusers/docs
|
||||
cp -r source/en source/LANG-ID
|
||||
```
|
||||
|
||||
Here, `LANG-ID` should be one of the ISO 639-1 or ISO 639-2 language codes -- see [here](https://www.loc.gov/standards/iso639-2/php/code_list.php) for a handy table.
|
||||
|
||||
**✍️ Start translating**
|
||||
|
||||
The fun part comes - translating the text!
|
||||
|
||||
The first thing we recommend is translating the part of the `_toctree.yml` file that corresponds to your doc chapter. This file is used to render the table of contents on the website.
|
||||
|
||||
> 🙋 If the `_toctree.yml` file doesn't yet exist for your language, you can create one by copy-pasting from the English version and deleting the sections unrelated to your chapter. Just make sure it exists in the `docs/source/LANG-ID/` directory!
|
||||
|
||||
The fields you should add are `local` (with the name of the file containing the translation; e.g. `autoclass_tutorial`), and `title` (with the title of the doc in your language; e.g. `Load pretrained instances with an AutoClass`) -- as a reference, here is the `_toctree.yml` for [English](https://github.com/huggingface/diffusers/blob/main/docs/source/en/_toctree.yml):
|
||||
|
||||
```yaml
|
||||
- sections:
|
||||
- local: pipeline_tutorial # Do not change this! Use the same name for your .md file
|
||||
title: Pipelines for inference # Translate this!
|
||||
...
|
||||
title: Tutorials # Translate this!
|
||||
```
|
||||
|
||||
Once you have translated the `_toctree.yml` file, you can start translating the [MDX](https://mdxjs.com/) files associated with your docs chapter.
|
||||
|
||||
> 🙋 If you'd like others to help you with the translation, you should [open an issue](https://github.com/huggingface/diffusers/issues) and tag @patrickvonplaten.
|
||||
100
docs/source/_toctree.yml
Normal file
100
docs/source/_toctree.yml
Normal file
@@ -0,0 +1,100 @@
|
||||
- sections:
|
||||
- local: index
|
||||
title: "🧨 Diffusers"
|
||||
- local: quicktour
|
||||
title: "Quicktour"
|
||||
- local: installation
|
||||
title: "Installation"
|
||||
title: "Get started"
|
||||
- sections:
|
||||
- sections:
|
||||
- local: using-diffusers/loading
|
||||
title: "Loading Pipelines, Models, and Schedulers"
|
||||
- local: using-diffusers/configuration
|
||||
title: "Configuring Pipelines, Models, and Schedulers"
|
||||
- local: using-diffusers/custom_pipeline_overview
|
||||
title: "Loading and Adding Custom Pipelines"
|
||||
title: "Loading & Hub"
|
||||
- sections:
|
||||
- local: using-diffusers/unconditional_image_generation
|
||||
title: "Unconditional Image Generation"
|
||||
- local: using-diffusers/conditional_image_generation
|
||||
title: "Text-to-Image Generation"
|
||||
- local: using-diffusers/img2img
|
||||
title: "Text-Guided Image-to-Image"
|
||||
- local: using-diffusers/inpaint
|
||||
title: "Text-Guided Image-Inpainting"
|
||||
- local: using-diffusers/custom_pipeline_examples
|
||||
title: "Community Pipelines"
|
||||
- local: using-diffusers/contribute_pipeline
|
||||
title: "How to contribute a Pipeline"
|
||||
title: "Pipelines for Inference"
|
||||
title: "Using Diffusers"
|
||||
- sections:
|
||||
- local: optimization/fp16
|
||||
title: "Memory and Speed"
|
||||
- local: optimization/onnx
|
||||
title: "ONNX"
|
||||
- local: optimization/open_vino
|
||||
title: "OpenVINO"
|
||||
- local: optimization/mps
|
||||
title: "MPS"
|
||||
title: "Optimization/Special Hardware"
|
||||
- sections:
|
||||
- local: training/overview
|
||||
title: "Overview"
|
||||
- local: training/unconditional_training
|
||||
title: "Unconditional Image Generation"
|
||||
- local: training/text_inversion
|
||||
title: "Textual Inversion"
|
||||
- local: training/dreambooth
|
||||
title: "Dreambooth"
|
||||
- local: training/text2image
|
||||
title: "Text-to-image fine-tuning"
|
||||
title: "Training"
|
||||
- sections:
|
||||
- local: conceptual/stable_diffusion
|
||||
title: "Stable Diffusion"
|
||||
- local: conceptual/philosophy
|
||||
title: "Philosophy"
|
||||
- local: conceptual/contribution
|
||||
title: "How to contribute?"
|
||||
title: "Conceptual Guides"
|
||||
- sections:
|
||||
- sections:
|
||||
- local: api/models
|
||||
title: "Models"
|
||||
- local: api/schedulers
|
||||
title: "Schedulers"
|
||||
- local: api/diffusion_pipeline
|
||||
title: "Diffusion Pipeline"
|
||||
- local: api/logging
|
||||
title: "Logging"
|
||||
- local: api/configuration
|
||||
title: "Configuration"
|
||||
- local: api/outputs
|
||||
title: "Outputs"
|
||||
title: "Main Classes"
|
||||
- sections:
|
||||
- local: api/pipelines/overview
|
||||
title: "Overview"
|
||||
- local: api/pipelines/ddim
|
||||
title: "DDIM"
|
||||
- local: api/pipelines/ddpm
|
||||
title: "DDPM"
|
||||
- local: api/pipelines/latent_diffusion
|
||||
title: "Latent Diffusion"
|
||||
- local: api/pipelines/latent_diffusion_uncond
|
||||
title: "Unconditional Latent Diffusion"
|
||||
- local: api/pipelines/pndm
|
||||
title: "PNDM"
|
||||
- local: api/pipelines/score_sde_ve
|
||||
title: "Score SDE VE"
|
||||
- local: api/pipelines/stable_diffusion
|
||||
title: "Stable Diffusion"
|
||||
- local: api/pipelines/stochastic_karras_ve
|
||||
title: "Stochastic Karras VE"
|
||||
- local: api/pipelines/dance_diffusion
|
||||
title: "Dance Diffusion"
|
||||
title: "Pipelines"
|
||||
title: "API"
|
||||
@@ -15,9 +15,9 @@ specific language governing permissions and limitations under the License.
|
||||
In Diffusers, schedulers of type [`schedulers.scheduling_utils.SchedulerMixin`], and models of type [`ModelMixin`] inherit from [`ConfigMixin`] which conveniently takes care of storing all parameters that are
|
||||
passed to the respective `__init__` methods in a JSON-configuration file.
|
||||
|
||||
## ConfigMixin
|
||||
TODO(PVP) - add example and better info here
|
||||
|
||||
## ConfigMixin
|
||||
[[autodoc]] ConfigMixin
|
||||
- load_config
|
||||
- from_config
|
||||
- save_config
|
||||
@@ -30,18 +30,13 @@ Any pipeline object can be saved locally with [`~DiffusionPipeline.save_pretrain
|
||||
|
||||
## DiffusionPipeline
|
||||
[[autodoc]] DiffusionPipeline
|
||||
- all
|
||||
- __call__
|
||||
- device
|
||||
- from_pretrained
|
||||
- save_pretrained
|
||||
- to
|
||||
- device
|
||||
- components
|
||||
|
||||
## ImagePipelineOutput
|
||||
By default diffusion pipelines return an object of class
|
||||
|
||||
[[autodoc]] pipelines.ImagePipelineOutput
|
||||
|
||||
## AudioPipelineOutput
|
||||
By default diffusion pipelines return an object of class
|
||||
|
||||
[[autodoc]] pipelines.AudioPipelineOutput
|
||||
[[autodoc]] pipeline_utils.ImagePipelineOutput
|
||||
@@ -1,4 +1,4 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
<!--Copyright 2020 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
@@ -22,15 +22,12 @@ The models are built on the base class ['ModelMixin'] that is a `torch.nn.module
|
||||
## UNet2DOutput
|
||||
[[autodoc]] models.unet_2d.UNet2DOutput
|
||||
|
||||
## UNet2DModel
|
||||
[[autodoc]] UNet2DModel
|
||||
|
||||
## UNet1DOutput
|
||||
[[autodoc]] models.unet_1d.UNet1DOutput
|
||||
|
||||
## UNet1DModel
|
||||
[[autodoc]] UNet1DModel
|
||||
|
||||
## UNet2DModel
|
||||
[[autodoc]] UNet2DModel
|
||||
|
||||
## UNet2DConditionOutput
|
||||
[[autodoc]] models.unet_2d_condition.UNet2DConditionOutput
|
||||
|
||||
@@ -41,29 +38,17 @@ The models are built on the base class ['ModelMixin'] that is a `torch.nn.module
|
||||
[[autodoc]] models.vae.DecoderOutput
|
||||
|
||||
## VQEncoderOutput
|
||||
[[autodoc]] models.vq_model.VQEncoderOutput
|
||||
[[autodoc]] models.vae.VQEncoderOutput
|
||||
|
||||
## VQModel
|
||||
[[autodoc]] VQModel
|
||||
|
||||
## AutoencoderKLOutput
|
||||
[[autodoc]] models.autoencoder_kl.AutoencoderKLOutput
|
||||
[[autodoc]] models.vae.AutoencoderKLOutput
|
||||
|
||||
## AutoencoderKL
|
||||
[[autodoc]] AutoencoderKL
|
||||
|
||||
## Transformer2DModel
|
||||
[[autodoc]] Transformer2DModel
|
||||
|
||||
## Transformer2DModelOutput
|
||||
[[autodoc]] models.transformer_2d.Transformer2DModelOutput
|
||||
|
||||
## PriorTransformer
|
||||
[[autodoc]] models.prior_transformer.PriorTransformer
|
||||
|
||||
## PriorTransformerOutput
|
||||
[[autodoc]] models.prior_transformer.PriorTransformerOutput
|
||||
|
||||
## FlaxModelMixin
|
||||
[[autodoc]] FlaxModelMixin
|
||||
|
||||
@@ -25,7 +25,7 @@ pipeline = DDIMPipeline.from_pretrained("google/ddpm-cifar10-32")
|
||||
outputs = pipeline()
|
||||
```
|
||||
|
||||
The `outputs` object is a [`~pipelines.ImagePipelineOutput`], as we can see in the
|
||||
The `outputs` object is a [`~pipeline_utils.ImagePipelineOutput`], as we can see in the
|
||||
documentation of that class below, it means it has an image attribute.
|
||||
|
||||
You can access each attribute as you would usually do, and if that attribute has not been returned by the model, you will get `None`:
|
||||
@@ -30,5 +30,4 @@ The original codebase of this implementation can be found [here](https://github.
|
||||
|
||||
## DanceDiffusionPipeline
|
||||
[[autodoc]] DanceDiffusionPipeline
|
||||
- all
|
||||
- __call__
|
||||
- __call__
|
||||
@@ -20,8 +20,7 @@ The abstract of the paper is the following:
|
||||
|
||||
Denoising diffusion probabilistic models (DDPMs) have achieved high quality image generation without adversarial training, yet they require simulating a Markov chain for many steps to produce a sample. To accelerate sampling, we present denoising diffusion implicit models (DDIMs), a more efficient class of iterative implicit probabilistic models with the same training procedure as DDPMs. In DDPMs, the generative process is defined as the reverse of a Markovian diffusion process. We construct a class of non-Markovian diffusion processes that lead to the same training objective, but whose reverse process can be much faster to sample from. We empirically demonstrate that DDIMs can produce high quality samples 10× to 50× faster in terms of wall-clock time compared to DDPMs, allow us to trade off computation for sample quality, and can perform semantically meaningful image interpolation directly in the latent space.
|
||||
|
||||
The original codebase of this paper can be found here: [ermongroup/ddim](https://github.com/ermongroup/ddim).
|
||||
For questions, feel free to contact the author on [tsong.me](https://tsong.me/).
|
||||
The original codebase of this paper can be found [here](https://github.com/ermongroup/ddim).
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
@@ -32,5 +31,4 @@ For questions, feel free to contact the author on [tsong.me](https://tsong.me/).
|
||||
|
||||
## DDIMPipeline
|
||||
[[autodoc]] DDIMPipeline
|
||||
- all
|
||||
- __call__
|
||||
- __call__
|
||||
@@ -33,5 +33,4 @@ The original codebase of this paper can be found [here](https://github.com/hojon
|
||||
|
||||
# DDPMPipeline
|
||||
[[autodoc]] DDPMPipeline
|
||||
- all
|
||||
- __call__
|
||||
- __call__
|
||||
@@ -33,17 +33,10 @@ The original codebase can be found [here](https://github.com/CompVis/latent-diff
|
||||
| Pipeline | Tasks | Colab
|
||||
|---|---|:---:|
|
||||
| [pipeline_latent_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/latent_diffusion/pipeline_latent_diffusion.py) | *Text-to-Image Generation* | - |
|
||||
| [pipeline_latent_diffusion_superresolution.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/latent_diffusion/pipeline_latent_diffusion_superresolution.py) | *Super Resolution* | - |
|
||||
|
||||
## Examples:
|
||||
|
||||
|
||||
## LDMTextToImagePipeline
|
||||
[[autodoc]] LDMTextToImagePipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## LDMSuperResolutionPipeline
|
||||
[[autodoc]] LDMSuperResolutionPipeline
|
||||
- all
|
||||
- __call__
|
||||
[[autodoc]] pipelines.latent_diffusion.pipeline_latent_diffusion.LDMTextToImagePipeline
|
||||
- __call__
|
||||
@@ -38,5 +38,4 @@ The original codebase can be found [here](https://github.com/CompVis/latent-diff
|
||||
|
||||
## LDMPipeline
|
||||
[[autodoc]] LDMPipeline
|
||||
- all
|
||||
- __call__
|
||||
- __call__
|
||||
@@ -28,7 +28,7 @@ or created independently from each other.
|
||||
|
||||
To that end, we strive to offer all open-sourced, state-of-the-art diffusion system under a unified API.
|
||||
More specifically, we strive to provide pipelines that
|
||||
- 1. can load the officially published weights and yield 1-to-1 the same outputs as the original implementation according to the corresponding paper (*e.g.* [LDMTextToImagePipeline](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines/latent_diffusion), uses the officially released weights of [High-Resolution Image Synthesis with Latent Diffusion Models](https://arxiv.org/abs/2112.10752)),
|
||||
- 1. can load the officially published weights and yield 1-to-1 the same outputs as the original implementation according to the corresponding paper (*e.g.* [LatentDiffusionPipeline](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines/latent_diffusion), uses the officially released weights of [High-Resolution Image Synthesis with Latent Diffusion Models](https://arxiv.org/abs/2112.10752)),
|
||||
- 2. have a simple user interface to run the model in inference (see the [Pipelines API](#pipelines-api) section),
|
||||
- 3. are easy to understand with code that is self-explanatory and can be read along-side the official paper (see [Pipelines summary](#pipelines-summary)),
|
||||
- 4. can easily be contributed by the community (see the [Contribution](#contribution) section).
|
||||
@@ -41,47 +41,19 @@ If you are looking for *official* training examples, please have a look at [exam
|
||||
The following table summarizes all officially supported pipelines, their corresponding paper, and if
|
||||
available a colab notebook to directly try them out.
|
||||
|
||||
|
||||
| Pipeline | Paper | Tasks | Colab
|
||||
|---|---|:---:|:---:|
|
||||
| [alt_diffusion](./alt_diffusion) | [**AltDiffusion**](https://arxiv.org/abs/2211.06679) | Image-to-Image Text-Guided Generation | -
|
||||
| [audio_diffusion](./audio_diffusion) | [**Audio Diffusion**](https://github.com/teticio/audio_diffusion.git) | Unconditional Audio Generation |
|
||||
| [cycle_diffusion](./cycle_diffusion) | [**Cycle Diffusion**](https://arxiv.org/abs/2210.05559) | Image-to-Image Text-Guided Generation |
|
||||
| [dance_diffusion](./dance_diffusion) | [**Dance Diffusion**](https://github.com/williamberman/diffusers.git) | Unconditional Audio Generation |
|
||||
| [ddpm](./ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) | Unconditional Image Generation |
|
||||
| [ddim](./ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) | Unconditional Image Generation |
|
||||
| [ddim](./ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) | Unconditional Image Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
|
||||
| [latent_diffusion](./latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation |
|
||||
| [latent_diffusion](./latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image |
|
||||
| [latent_diffusion_uncond](./latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation |
|
||||
| [paint_by_example](./paint_by_example) | [**Paint by Example: Exemplar-based Image Editing with Diffusion Models**](https://arxiv.org/abs/2211.13227) | Image-Guided Image Inpainting |
|
||||
| [pndm](./pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation |
|
||||
| [score_sde_ve](./score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
|
||||
| [score_sde_vp](./score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
|
||||
| [semantic_stable_diffusion](./semantic_stable_diffusion) | [**SEGA: Instructing Diffusion using Semantic Dimensions**](https://arxiv.org/abs/2301.12247) | Text-to-Image Generation |
|
||||
| [stable_diffusion_text2img](./stable_diffusion/text2img) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
|
||||
| [stable_diffusion_img2img](./stable_diffusion/img2img) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
|
||||
| [stable_diffusion_inpaint](./stable_diffusion/inpaint) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
|
||||
| [stable_diffusion_panorama](./stable_diffusion/panorama) | [**MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation**](https://arxiv.org/abs/2302.08113) | Text-Guided Panorama View Generation |
|
||||
| [stable_diffusion_pix2pix](./stable_diffusion/pix2pix) | [**InstructPix2Pix: Learning to Follow Image Editing Instructions**](https://arxiv.org/abs/2211.09800) | Text-Based Image Editing |
|
||||
| [stable_diffusion_pix2pix_zero](./stable_diffusion/pix2pix_zero) | [**Zero-shot Image-to-Image Translation**](https://arxiv.org/abs/2302.03027) | Text-Based Image Editing |
|
||||
| [stable_diffusion_attend_and_excite](./stable_diffusion/attend_and_excite) | [**Attend and Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models**](https://arxiv.org/abs/2301.13826) | Text-to-Image Generation |
|
||||
| [stable_diffusion_self_attention_guidance](./stable_diffusion/self_attention_guidance) | [**Self-Attention Guidance**](https://arxiv.org/abs/2210.00939) | Text-to-Image Generation |
|
||||
| [stable_diffusion_image_variation](./stable_diffusion/image_variation) | [**Stable Diffusion Image Variations**](https://github.com/LambdaLabsML/lambda-diffusers#stable-diffusion-image-variations) | Image-to-Image Generation |
|
||||
| [stable_diffusion_latent_upscale](./stable_diffusion/latent_upscale) | [**Stable Diffusion Latent Upscaler**](https://twitter.com/StabilityAI/status/1590531958815064065) | Text-Guided Super Resolution Image-to-Image |
|
||||
| [stable_diffusion_2](./stable_diffusion_2/) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-to-Image Generation |
|
||||
| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting |
|
||||
| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Depth-to-Image Text-Guided Generation |
|
||||
| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
|
||||
| [stable_diffusion_safe](./stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
|
||||
| [stable_unclip](./stable_unclip) | **Stable unCLIP** | Text-to-Image Generation |
|
||||
| [stable_unclip](./stable_unclip) | **Stable unCLIP** | Image-to-Image Text-Guided Generation |
|
||||
| [stochastic_karras_ve](./stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
|
||||
| [unclip](./unclip) | [Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) | Text-to-Image Generation |
|
||||
| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
|
||||
| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation |
|
||||
| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |
|
||||
| [vq_diffusion](./vq_diffusion) | [Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation |
|
||||
|
||||
| [stable_diffusion](./stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
|
||||
| [stable_diffusion](./stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
|
||||
| [stable_diffusion](./stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
|
||||
| [stochastic_karras_ve](./stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
|
||||
|
||||
**Note**: Pipelines are simple examples of how to play around with the diffusion systems as described in the corresponding papers.
|
||||
|
||||
@@ -150,9 +122,9 @@ from diffusers import StableDiffusionImg2ImgPipeline
|
||||
|
||||
# load the pipeline
|
||||
device = "cuda"
|
||||
pipe = StableDiffusionImg2ImgPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16).to(
|
||||
device
|
||||
)
|
||||
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
|
||||
"runwayml/stable-diffusion-v1-5", revision="fp16", torch_dtype=torch.float16
|
||||
).to(device)
|
||||
|
||||
# let's download an initial image
|
||||
url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
|
||||
@@ -163,7 +135,7 @@ init_image = init_image.resize((768, 512))
|
||||
|
||||
prompt = "A fantasy landscape, trending on artstation"
|
||||
|
||||
images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
|
||||
images[0].save("fantasy_landscape.png")
|
||||
```
|
||||
@@ -200,6 +172,7 @@ mask_image = download_image(mask_url).resize((512, 512))
|
||||
|
||||
pipe = StableDiffusionInpaintPipeline.from_pretrained(
|
||||
"runwayml/stable-diffusion-inpainting",
|
||||
revision="fp16",
|
||||
torch_dtype=torch.float16,
|
||||
)
|
||||
pipe = pipe.to("cuda")
|
||||
@@ -30,6 +30,6 @@ The original codebase can be found [here](https://github.com/luping-liu/PNDM).
|
||||
|
||||
|
||||
## PNDMPipeline
|
||||
[[autodoc]] PNDMPipeline
|
||||
- all
|
||||
- __call__
|
||||
[[autodoc]] pipelines.pndm.pipeline_pndm.PNDMPipeline
|
||||
- __call__
|
||||
|
||||
@@ -32,5 +32,5 @@ This pipeline implements the Variance Expanding (VE) variant of the method.
|
||||
|
||||
## ScoreSdeVePipeline
|
||||
[[autodoc]] ScoreSdeVePipeline
|
||||
- all
|
||||
- __call__
|
||||
- __call__
|
||||
|
||||
71
docs/source/api/pipelines/stable_diffusion.mdx
Normal file
71
docs/source/api/pipelines/stable_diffusion.mdx
Normal file
@@ -0,0 +1,71 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Stable diffusion pipelines
|
||||
|
||||
Stable Diffusion is a text-to-image _latent diffusion_ model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/). It's trained on 512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/) dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs.
|
||||
|
||||
Latent diffusion is the research on top of which Stable Diffusion was built. It was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://arxiv.org/abs/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer. You can learn more details about it in the [specific pipeline for latent diffusion](pipelines/latent_diffusion) that is part of 🤗 Diffusers.
|
||||
|
||||
For more details about how Stable Diffusion works and how it differs from the base latent diffusion model, please refer to the official [launch announcement post](https://stability.ai/blog/stable-diffusion-announcement) and [this section of our own blog post](https://huggingface.co/blog/stable_diffusion#how-does-stable-diffusion-work).
|
||||
|
||||
*Tips*:
|
||||
- To tweak your prompts on a specific result you liked, you can generate your own latents, as demonstrated in the following notebook: [](https://colab.research.google.com/github/pcuenca/diffusers-examples/blob/main/notebooks/stable-diffusion-seeds.ipynb)
|
||||
|
||||
*Overview*:
|
||||
|
||||
| Pipeline | Tasks | Colab | Demo
|
||||
|---|---|:---:|:---:|
|
||||
| [pipeline_stable_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py) | *Text-to-Image Generation* | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb) | [🤗 Stable Diffusion](https://huggingface.co/spaces/stabilityai/stable-diffusion)
|
||||
| [pipeline_stable_diffusion_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_img2img.py) | *Image-to-Image Text-Guided Generation* | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb) | [🤗 Diffuse the Rest](https://huggingface.co/spaces/huggingface/diffuse-the-rest)
|
||||
| [pipeline_stable_diffusion_inpaint.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_inpaint.py) | **Experimental** – *Text-Guided Image Inpainting* | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb) | Coming soon
|
||||
|
||||
## Tips
|
||||
|
||||
If you want to use all possible use cases in a single `DiffusionPipeline` you can either:
|
||||
- Make use of the [Stable Diffusion Mega Pipeline](https://github.com/huggingface/diffusers/tree/main/examples/community#stable-diffusion-mega) or
|
||||
- Make use of the `components` functionality to instantiate all components in the most memory-efficient way:
|
||||
|
||||
```python
|
||||
>>> from diffusers import (
|
||||
... StableDiffusionPipeline,
|
||||
... StableDiffusionImg2ImgPipeline,
|
||||
... StableDiffusionInpaintPipeline,
|
||||
... )
|
||||
|
||||
>>> img2text = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4")
|
||||
>>> img2img = StableDiffusionImg2ImgPipeline(**img2text.components)
|
||||
>>> inpaint = StableDiffusionInpaintPipeline(**img2text.components)
|
||||
|
||||
>>> # now you can use img2text(...), img2img(...), inpaint(...) just like the call methods of each respective pipeline
|
||||
```
|
||||
|
||||
## StableDiffusionPipelineOutput
|
||||
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
|
||||
|
||||
## StableDiffusionPipeline
|
||||
[[autodoc]] StableDiffusionPipeline
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
|
||||
## StableDiffusionImg2ImgPipeline
|
||||
[[autodoc]] StableDiffusionImg2ImgPipeline
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
|
||||
## StableDiffusionInpaintPipeline
|
||||
[[autodoc]] StableDiffusionInpaintPipeline
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
@@ -32,5 +32,4 @@ This pipeline implements the Stochastic sampling tailored to the Variance-Expand
|
||||
|
||||
## KarrasVePipeline
|
||||
[[autodoc]] KarrasVePipeline
|
||||
- all
|
||||
- __call__
|
||||
- __call__
|
||||
114
docs/source/api/schedulers.mdx
Normal file
114
docs/source/api/schedulers.mdx
Normal file
@@ -0,0 +1,114 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Schedulers
|
||||
|
||||
Diffusers contains multiple pre-built schedule functions for the diffusion process.
|
||||
|
||||
## What is a scheduler?
|
||||
|
||||
The schedule functions, denoted *Schedulers* in the library take in the output of a trained model, a sample which the diffusion process is iterating on, and a timestep to return a denoised sample.
|
||||
|
||||
- Schedulers define the methodology for iteratively adding noise to an image or for updating a sample based on model outputs.
|
||||
- adding noise in different manners represent the algorithmic processes to train a diffusion model by adding noise to images.
|
||||
- for inference, the scheduler defines how to update a sample based on an output from a pretrained model.
|
||||
- Schedulers are often defined by a *noise schedule* and an *update rule* to solve the differential equation solution.
|
||||
|
||||
### Discrete versus continuous schedulers
|
||||
|
||||
All schedulers take in a timestep to predict the updated version of the sample being diffused.
|
||||
The timesteps dictate where in the diffusion process the step is, where data is generated by iterating forward in time and inference is executed by propagating backwards through timesteps.
|
||||
Different algorithms use timesteps that both discrete (accepting `int` inputs), such as the [`DDPMScheduler`] or [`PNDMScheduler`], and continuous (accepting `float` inputs), such as the score-based schedulers [`ScoreSdeVeScheduler`] or [`ScoreSdeVpScheduler`].
|
||||
|
||||
## Designing Re-usable schedulers
|
||||
|
||||
The core design principle between the schedule functions is to be model, system, and framework independent.
|
||||
This allows for rapid experimentation and cleaner abstractions in the code, where the model prediction is separated from the sample update.
|
||||
To this end, the design of schedulers is such that:
|
||||
|
||||
- Schedulers can be used interchangeably between diffusion models in inference to find the preferred trade-off between speed and generation quality.
|
||||
- Schedulers are currently by default in PyTorch, but are designed to be framework independent (partial Jax support currently exists).
|
||||
|
||||
|
||||
## API
|
||||
|
||||
The core API for any new scheduler must follow a limited structure.
|
||||
- Schedulers should provide one or more `def step(...)` functions that should be called to update the generated sample iteratively.
|
||||
- Schedulers should provide a `set_timesteps(...)` method that configures the parameters of a schedule function for a specific inference task.
|
||||
- Schedulers should be framework-specific.
|
||||
|
||||
The base class [`SchedulerMixin`] implements low level utilities used by multiple schedulers.
|
||||
|
||||
### SchedulerMixin
|
||||
[[autodoc]] SchedulerMixin
|
||||
|
||||
### SchedulerOutput
|
||||
The class [`SchedulerOutput`] contains the outputs from any schedulers `step(...)` call.
|
||||
|
||||
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
|
||||
|
||||
### Implemented Schedulers
|
||||
|
||||
#### Denoising diffusion implicit models (DDIM)
|
||||
|
||||
Original paper can be found here.
|
||||
|
||||
[[autodoc]] DDIMScheduler
|
||||
|
||||
#### Denoising diffusion probabilistic models (DDPM)
|
||||
|
||||
Original paper can be found [here](https://arxiv.org/abs/2010.02502).
|
||||
|
||||
[[autodoc]] DDPMScheduler
|
||||
|
||||
#### Variance exploding, stochastic sampling from Karras et. al
|
||||
|
||||
Original paper can be found [here](https://arxiv.org/abs/2006.11239).
|
||||
|
||||
[[autodoc]] KarrasVeScheduler
|
||||
|
||||
#### Linear multistep scheduler for discrete beta schedules
|
||||
|
||||
Original implementation can be found [here](https://arxiv.org/abs/2206.00364).
|
||||
|
||||
|
||||
[[autodoc]] LMSDiscreteScheduler
|
||||
|
||||
#### Pseudo numerical methods for diffusion models (PNDM)
|
||||
|
||||
Original implementation can be found [here](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181).
|
||||
|
||||
[[autodoc]] PNDMScheduler
|
||||
|
||||
#### variance exploding stochastic differential equation (VE-SDE) scheduler
|
||||
|
||||
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
|
||||
|
||||
[[autodoc]] ScoreSdeVeScheduler
|
||||
|
||||
#### improved pseudo numerical methods for diffusion models (iPNDM)
|
||||
|
||||
Original implementation can be found [here](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296).
|
||||
|
||||
[[autodoc]] IPNDMScheduler
|
||||
|
||||
#### variance preserving stochastic differential equation (VP-SDE) scheduler
|
||||
|
||||
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
|
||||
|
||||
<Tip warning={true}>
|
||||
|
||||
Score SDE-VP is under construction.
|
||||
|
||||
</Tip>
|
||||
|
||||
[[autodoc]] schedulers.scheduling_sde_vp.ScoreSdeVpScheduler
|
||||
@@ -177,7 +177,7 @@ Follow these steps to start contributing ([supported Python versions](https://gi
|
||||
$ make style
|
||||
```
|
||||
|
||||
🧨 Diffusers also uses `ruff` and a few custom scripts to check for coding mistakes. Quality
|
||||
🧨 Diffusers also uses `flake8` and a few custom scripts to check for coding mistakes. Quality
|
||||
control runs in CI, however you can also run the same checks with:
|
||||
|
||||
```bash
|
||||
17
docs/source/conceptual/philosophy.mdx
Normal file
17
docs/source/conceptual/philosophy.mdx
Normal file
@@ -0,0 +1,17 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Philosophy
|
||||
|
||||
- Readability and clarity are preferred over highly optimized code. A strong importance is put on providing readable, intuitive and elementary code design. *E.g.*, the provided [schedulers](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers) are separated from the provided [models](https://github.com/huggingface/diffusers/tree/main/src/diffusers/models) and use well-commented code that can be read alongside the original paper.
|
||||
- Diffusers is **modality independent** and focuses on providing pretrained models and tools to build systems that generate **continuous outputs**, *e.g.* vision and audio. This is one of the guiding goals even if the initial pipelines are devoted to vision tasks.
|
||||
- Diffusion models and schedulers are provided as concise, elementary building blocks. In contrast, diffusion pipelines are a collection of end-to-end diffusion systems that can be used out-of-the-box, should stay as close as possible to their original implementations and can include components of other libraries, such as text encoders. Examples of diffusion pipelines are [Glide](https://github.com/openai/glide-text2im), [Latent Diffusion](https://github.com/CompVis/latent-diffusion) and [Stable Diffusion](https://github.com/compvis/stable-diffusion).
|
||||
@@ -10,7 +10,6 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# 번역중
|
||||
# Stable Diffusion
|
||||
|
||||
열심히 번역을 진행중입니다. 조금만 기다려주세요.
|
||||
감사합니다!
|
||||
Please visit this [very in-detail blog post](https://huggingface.co/blog/stable_diffusion) on Stable Diffusion!
|
||||
@@ -1,232 +0,0 @@
|
||||
- sections:
|
||||
- local: index
|
||||
title: 🧨 Diffusers
|
||||
- local: quicktour
|
||||
title: Quicktour
|
||||
- local: stable_diffusion
|
||||
title: Stable Diffusion
|
||||
- local: installation
|
||||
title: Installation
|
||||
title: Get started
|
||||
- sections:
|
||||
- sections:
|
||||
- local: using-diffusers/loading
|
||||
title: Loading Pipelines, Models, and Schedulers
|
||||
- local: using-diffusers/schedulers
|
||||
title: Using different Schedulers
|
||||
- local: using-diffusers/configuration
|
||||
title: Configuring Pipelines, Models, and Schedulers
|
||||
- local: using-diffusers/custom_pipeline_overview
|
||||
title: Loading and Adding Custom Pipelines
|
||||
- local: using-diffusers/kerascv
|
||||
title: Using KerasCV Stable Diffusion Checkpoints in Diffusers
|
||||
title: Loading & Hub
|
||||
- sections:
|
||||
- local: using-diffusers/unconditional_image_generation
|
||||
title: Unconditional Image Generation
|
||||
- local: using-diffusers/conditional_image_generation
|
||||
title: Text-to-Image Generation
|
||||
- local: using-diffusers/img2img
|
||||
title: Text-Guided Image-to-Image
|
||||
- local: using-diffusers/inpaint
|
||||
title: Text-Guided Image-Inpainting
|
||||
- local: using-diffusers/depth2img
|
||||
title: Text-Guided Depth-to-Image
|
||||
- local: using-diffusers/controlling_generation
|
||||
title: Controlling generation
|
||||
- local: using-diffusers/reusing_seeds
|
||||
title: Reusing seeds for deterministic generation
|
||||
- local: using-diffusers/reproducibility
|
||||
title: Reproducibility
|
||||
- local: using-diffusers/custom_pipeline_examples
|
||||
title: Community Pipelines
|
||||
- local: using-diffusers/contribute_pipeline
|
||||
title: How to contribute a Pipeline
|
||||
- local: using-diffusers/using_safetensors
|
||||
title: Using safetensors
|
||||
title: Pipelines for Inference
|
||||
- sections:
|
||||
- local: using-diffusers/rl
|
||||
title: Reinforcement Learning
|
||||
- local: using-diffusers/audio
|
||||
title: Audio
|
||||
- local: using-diffusers/other-modalities
|
||||
title: Other Modalities
|
||||
title: Taking Diffusers Beyond Images
|
||||
title: Using Diffusers
|
||||
- sections:
|
||||
- local: optimization/fp16
|
||||
title: Memory and Speed
|
||||
- local: optimization/torch2.0
|
||||
title: Torch2.0 support
|
||||
- local: optimization/xformers
|
||||
title: xFormers
|
||||
- local: optimization/onnx
|
||||
title: ONNX
|
||||
- local: optimization/open_vino
|
||||
title: OpenVINO
|
||||
- local: optimization/mps
|
||||
title: MPS
|
||||
- local: optimization/habana
|
||||
title: Habana Gaudi
|
||||
title: Optimization/Special Hardware
|
||||
- sections:
|
||||
- local: training/overview
|
||||
title: Overview
|
||||
- local: training/unconditional_training
|
||||
title: Unconditional Image Generation
|
||||
- local: training/text_inversion
|
||||
title: Textual Inversion
|
||||
- local: training/dreambooth
|
||||
title: Dreambooth
|
||||
- local: training/text2image
|
||||
title: Text-to-image fine-tuning
|
||||
- local: training/lora
|
||||
title: LoRA Support in Diffusers
|
||||
title: Training
|
||||
- sections:
|
||||
- local: conceptual/philosophy
|
||||
title: Philosophy
|
||||
- local: conceptual/contribution
|
||||
title: How to contribute?
|
||||
- local: conceptual/ethical_guidelines
|
||||
title: Diffusers' Ethical Guidelines
|
||||
title: Conceptual Guides
|
||||
- sections:
|
||||
- sections:
|
||||
- local: api/models
|
||||
title: Models
|
||||
- local: api/diffusion_pipeline
|
||||
title: Diffusion Pipeline
|
||||
- local: api/logging
|
||||
title: Logging
|
||||
- local: api/configuration
|
||||
title: Configuration
|
||||
- local: api/outputs
|
||||
title: Outputs
|
||||
- local: api/loaders
|
||||
title: Loaders
|
||||
title: Main Classes
|
||||
- sections:
|
||||
- local: api/pipelines/overview
|
||||
title: Overview
|
||||
- local: api/pipelines/alt_diffusion
|
||||
title: AltDiffusion
|
||||
- local: api/pipelines/audio_diffusion
|
||||
title: Audio Diffusion
|
||||
- local: api/pipelines/cycle_diffusion
|
||||
title: Cycle Diffusion
|
||||
- local: api/pipelines/dance_diffusion
|
||||
title: Dance Diffusion
|
||||
- local: api/pipelines/ddim
|
||||
title: DDIM
|
||||
- local: api/pipelines/ddpm
|
||||
title: DDPM
|
||||
- local: api/pipelines/dit
|
||||
title: DiT
|
||||
- local: api/pipelines/latent_diffusion
|
||||
title: Latent Diffusion
|
||||
- local: api/pipelines/paint_by_example
|
||||
title: PaintByExample
|
||||
- local: api/pipelines/pndm
|
||||
title: PNDM
|
||||
- local: api/pipelines/repaint
|
||||
title: RePaint
|
||||
- local: api/pipelines/stable_diffusion_safe
|
||||
title: Safe Stable Diffusion
|
||||
- local: api/pipelines/score_sde_ve
|
||||
title: Score SDE VE
|
||||
- local: api/pipelines/semantic_stable_diffusion
|
||||
title: Semantic Guidance
|
||||
- sections:
|
||||
- local: api/pipelines/stable_diffusion/overview
|
||||
title: Overview
|
||||
- local: api/pipelines/stable_diffusion/text2img
|
||||
title: Text-to-Image
|
||||
- local: api/pipelines/stable_diffusion/img2img
|
||||
title: Image-to-Image
|
||||
- local: api/pipelines/stable_diffusion/inpaint
|
||||
title: Inpaint
|
||||
- local: api/pipelines/stable_diffusion/depth2img
|
||||
title: Depth-to-Image
|
||||
- local: api/pipelines/stable_diffusion/image_variation
|
||||
title: Image-Variation
|
||||
- local: api/pipelines/stable_diffusion/upscale
|
||||
title: Super-Resolution
|
||||
- local: api/pipelines/stable_diffusion/latent_upscale
|
||||
title: Stable-Diffusion-Latent-Upscaler
|
||||
- local: api/pipelines/stable_diffusion/pix2pix
|
||||
title: InstructPix2Pix
|
||||
- local: api/pipelines/stable_diffusion/attend_and_excite
|
||||
title: Attend and Excite
|
||||
- local: api/pipelines/stable_diffusion/pix2pix_zero
|
||||
title: Pix2Pix Zero
|
||||
- local: api/pipelines/stable_diffusion/self_attention_guidance
|
||||
title: Self-Attention Guidance
|
||||
- local: api/pipelines/stable_diffusion/panorama
|
||||
title: MultiDiffusion Panorama
|
||||
title: Stable Diffusion
|
||||
- local: api/pipelines/stable_diffusion_2
|
||||
title: Stable Diffusion 2
|
||||
- local: api/pipelines/stable_unclip
|
||||
title: Stable unCLIP
|
||||
- local: api/pipelines/stochastic_karras_ve
|
||||
title: Stochastic Karras VE
|
||||
- local: api/pipelines/unclip
|
||||
title: UnCLIP
|
||||
- local: api/pipelines/latent_diffusion_uncond
|
||||
title: Unconditional Latent Diffusion
|
||||
- local: api/pipelines/versatile_diffusion
|
||||
title: Versatile Diffusion
|
||||
- local: api/pipelines/vq_diffusion
|
||||
title: VQ Diffusion
|
||||
title: Pipelines
|
||||
- sections:
|
||||
- local: api/schedulers/overview
|
||||
title: Overview
|
||||
- local: api/schedulers/ddim
|
||||
title: DDIM
|
||||
- local: api/schedulers/ddim_inverse
|
||||
title: DDIMInverse
|
||||
- local: api/schedulers/ddpm
|
||||
title: DDPM
|
||||
- local: api/schedulers/deis
|
||||
title: DEIS
|
||||
- local: api/schedulers/dpm_discrete
|
||||
title: DPM Discrete Scheduler
|
||||
- local: api/schedulers/dpm_discrete_ancestral
|
||||
title: DPM Discrete Scheduler with ancestral sampling
|
||||
- local: api/schedulers/euler_ancestral
|
||||
title: Euler Ancestral Scheduler
|
||||
- local: api/schedulers/euler
|
||||
title: Euler scheduler
|
||||
- local: api/schedulers/heun
|
||||
title: Heun Scheduler
|
||||
- local: api/schedulers/ipndm
|
||||
title: IPNDM
|
||||
- local: api/schedulers/lms_discrete
|
||||
title: Linear Multistep
|
||||
- local: api/schedulers/multistep_dpm_solver
|
||||
title: Multistep DPM-Solver
|
||||
- local: api/schedulers/pndm
|
||||
title: PNDM
|
||||
- local: api/schedulers/repaint
|
||||
title: RePaint Scheduler
|
||||
- local: api/schedulers/singlestep_dpm_solver
|
||||
title: Singlestep DPM-Solver
|
||||
- local: api/schedulers/stochastic_karras_ve
|
||||
title: Stochastic Kerras VE
|
||||
- local: api/schedulers/unipc
|
||||
title: UniPCMultistepScheduler
|
||||
- local: api/schedulers/score_sde_ve
|
||||
title: VE-SDE
|
||||
- local: api/schedulers/score_sde_vp
|
||||
title: VP-SDE
|
||||
- local: api/schedulers/vq_diffusion
|
||||
title: VQDiffusionScheduler
|
||||
title: Schedulers
|
||||
- sections:
|
||||
- local: api/experimental/rl
|
||||
title: RL Planning
|
||||
title: Experimental Features
|
||||
title: API
|
||||
@@ -1,15 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# TODO
|
||||
|
||||
Coming soon!
|
||||
@@ -1,30 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Loaders
|
||||
|
||||
There are many ways to train adapter neural networks for diffusion models, such as
|
||||
- [Textual Inversion](./training/text_inversion.mdx)
|
||||
- [LoRA](https://github.com/cloneofsimo/lora)
|
||||
- [Hypernetworks](https://arxiv.org/abs/1609.09106)
|
||||
|
||||
Such adapter neural networks often only consist of a fraction of the number of weights compared
|
||||
to the pretrained model and as such are very portable. The Diffusers library offers an easy-to-use
|
||||
API to load such adapter neural networks via the [`loaders.py` module](https://github.com/huggingface/diffusers/blob/main/src/diffusers/loaders.py).
|
||||
|
||||
**Note**: This module is still highly experimental and prone to future changes.
|
||||
|
||||
## LoaderMixins
|
||||
|
||||
### UNet2DConditionLoadersMixin
|
||||
|
||||
[[autodoc]] loaders.UNet2DConditionLoadersMixin
|
||||
@@ -1,83 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# AltDiffusion
|
||||
|
||||
AltDiffusion was proposed in [AltCLIP: Altering the Language Encoder in CLIP for Extended Language Capabilities](https://arxiv.org/abs/2211.06679) by Zhongzhi Chen, Guang Liu, Bo-Wen Zhang, Fulong Ye, Qinghong Yang, Ledell Wu
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*In this work, we present a conceptually simple and effective method to train a strong bilingual multimodal representation model. Starting from the pretrained multimodal representation model CLIP released by OpenAI, we switched its text encoder with a pretrained multilingual text encoder XLM-R, and aligned both languages and image representations by a two-stage training schema consisting of teacher learning and contrastive learning. We validate our method through evaluations of a wide range of tasks. We set new state-of-the-art performances on a bunch of tasks including ImageNet-CN, Flicker30k- CN, and COCO-CN. Further, we obtain very close performances with CLIP on almost all tasks, suggesting that one can simply alter the text encoder in CLIP for extended capabilities such as multilingual understanding.*
|
||||
|
||||
|
||||
*Overview*:
|
||||
|
||||
| Pipeline | Tasks | Colab | Demo
|
||||
|---|---|:---:|:---:|
|
||||
| [pipeline_alt_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/alt_diffusion/pipeline_alt_diffusion.py) | *Text-to-Image Generation* | - | -
|
||||
| [pipeline_alt_diffusion_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/alt_diffusion/pipeline_alt_diffusion_img2img.py) | *Image-to-Image Text-Guided Generation* | - |-
|
||||
|
||||
## Tips
|
||||
|
||||
- AltDiffusion is conceptually exaclty the same as [Stable Diffusion](./api/pipelines/stable_diffusion/overview).
|
||||
|
||||
- *Run AltDiffusion*
|
||||
|
||||
AltDiffusion can be tested very easily with the [`AltDiffusionPipeline`], [`AltDiffusionImg2ImgPipeline`] and the `"BAAI/AltDiffusion-m9"` checkpoint exactly in the same way it is shown in the [Conditional Image Generation Guide](./using-diffusers/conditional_image_generation) and the [Image-to-Image Generation Guide](./using-diffusers/img2img).
|
||||
|
||||
- *How to load and use different schedulers.*
|
||||
|
||||
The alt diffusion pipeline uses [`DDIMScheduler`] scheduler by default. But `diffusers` provides many other schedulers that can be used with the alt diffusion pipeline such as [`PNDMScheduler`], [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`] etc.
|
||||
To use a different scheduler, you can either change it via the [`ConfigMixin.from_config`] method or pass the `scheduler` argument to the `from_pretrained` method of the pipeline. For example, to use the [`EulerDiscreteScheduler`], you can do the following:
|
||||
|
||||
```python
|
||||
>>> from diffusers import AltDiffusionPipeline, EulerDiscreteScheduler
|
||||
|
||||
>>> pipeline = AltDiffusionPipeline.from_pretrained("BAAI/AltDiffusion-m9")
|
||||
>>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
|
||||
|
||||
>>> # or
|
||||
>>> euler_scheduler = EulerDiscreteScheduler.from_pretrained("BAAI/AltDiffusion-m9", subfolder="scheduler")
|
||||
>>> pipeline = AltDiffusionPipeline.from_pretrained("BAAI/AltDiffusion-m9", scheduler=euler_scheduler)
|
||||
```
|
||||
|
||||
|
||||
- *How to convert all use cases with multiple or single pipeline*
|
||||
|
||||
If you want to use all possible use cases in a single `DiffusionPipeline` we recommend using the `components` functionality to instantiate all components in the most memory-efficient way:
|
||||
|
||||
```python
|
||||
>>> from diffusers import (
|
||||
... AltDiffusionPipeline,
|
||||
... AltDiffusionImg2ImgPipeline,
|
||||
... )
|
||||
|
||||
>>> text2img = AltDiffusionPipeline.from_pretrained("BAAI/AltDiffusion-m9")
|
||||
>>> img2img = AltDiffusionImg2ImgPipeline(**text2img.components)
|
||||
|
||||
>>> # now you can use text2img(...) and img2img(...) just like the call methods of each respective pipeline
|
||||
```
|
||||
|
||||
## AltDiffusionPipelineOutput
|
||||
[[autodoc]] pipelines.alt_diffusion.AltDiffusionPipelineOutput
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## AltDiffusionPipeline
|
||||
[[autodoc]] AltDiffusionPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## AltDiffusionImg2ImgPipeline
|
||||
[[autodoc]] AltDiffusionImg2ImgPipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,98 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Audio Diffusion
|
||||
|
||||
## Overview
|
||||
|
||||
[Audio Diffusion](https://github.com/teticio/audio-diffusion) by Robert Dargavel Smith.
|
||||
|
||||
Audio Diffusion leverages the recent advances in image generation using diffusion models by converting audio samples to
|
||||
and from mel spectrogram images.
|
||||
|
||||
The original codebase of this implementation can be found [here](https://github.com/teticio/audio-diffusion), including
|
||||
training scripts and example notebooks.
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Colab
|
||||
|---|---|:---:|
|
||||
| [pipeline_audio_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/audio_diffusion/pipeline_audio_diffusion.py) | *Unconditional Audio Generation* | [](https://colab.research.google.com/github/teticio/audio-diffusion/blob/master/notebooks/audio_diffusion_pipeline.ipynb) |
|
||||
|
||||
|
||||
## Examples:
|
||||
|
||||
### Audio Diffusion
|
||||
|
||||
```python
|
||||
import torch
|
||||
from IPython.display import Audio
|
||||
from diffusers import DiffusionPipeline
|
||||
|
||||
device = "cuda" if torch.cuda.is_available() else "cpu"
|
||||
pipe = DiffusionPipeline.from_pretrained("teticio/audio-diffusion-256").to(device)
|
||||
|
||||
output = pipe()
|
||||
display(output.images[0])
|
||||
display(Audio(output.audios[0], rate=mel.get_sample_rate()))
|
||||
```
|
||||
|
||||
### Latent Audio Diffusion
|
||||
|
||||
```python
|
||||
import torch
|
||||
from IPython.display import Audio
|
||||
from diffusers import DiffusionPipeline
|
||||
|
||||
device = "cuda" if torch.cuda.is_available() else "cpu"
|
||||
pipe = DiffusionPipeline.from_pretrained("teticio/latent-audio-diffusion-256").to(device)
|
||||
|
||||
output = pipe()
|
||||
display(output.images[0])
|
||||
display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
|
||||
```
|
||||
|
||||
### Audio Diffusion with DDIM (faster)
|
||||
|
||||
```python
|
||||
import torch
|
||||
from IPython.display import Audio
|
||||
from diffusers import DiffusionPipeline
|
||||
|
||||
device = "cuda" if torch.cuda.is_available() else "cpu"
|
||||
pipe = DiffusionPipeline.from_pretrained("teticio/audio-diffusion-ddim-256").to(device)
|
||||
|
||||
output = pipe()
|
||||
display(output.images[0])
|
||||
display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
|
||||
```
|
||||
|
||||
### Variations, in-painting, out-painting etc.
|
||||
|
||||
```python
|
||||
output = pipe(
|
||||
raw_audio=output.audios[0, 0],
|
||||
start_step=int(pipe.get_default_steps() / 2),
|
||||
mask_start_secs=1,
|
||||
mask_end_secs=1,
|
||||
)
|
||||
display(output.images[0])
|
||||
display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
|
||||
```
|
||||
|
||||
## AudioDiffusionPipeline
|
||||
[[autodoc]] AudioDiffusionPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## Mel
|
||||
[[autodoc]] Mel
|
||||
@@ -1,100 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Cycle Diffusion
|
||||
|
||||
## Overview
|
||||
|
||||
Cycle Diffusion is a Text-Guided Image-to-Image Generation model proposed in [Unifying Diffusion Models' Latent Space, with Applications to CycleDiffusion and Guidance](https://arxiv.org/abs/2210.05559) by Chen Henry Wu, Fernando De la Torre.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*Diffusion models have achieved unprecedented performance in generative modeling. The commonly-adopted formulation of the latent code of diffusion models is a sequence of gradually denoised samples, as opposed to the simpler (e.g., Gaussian) latent space of GANs, VAEs, and normalizing flows. This paper provides an alternative, Gaussian formulation of the latent space of various diffusion models, as well as an invertible DPM-Encoder that maps images into the latent space. While our formulation is purely based on the definition of diffusion models, we demonstrate several intriguing consequences. (1) Empirically, we observe that a common latent space emerges from two diffusion models trained independently on related domains. In light of this finding, we propose CycleDiffusion, which uses DPM-Encoder for unpaired image-to-image translation. Furthermore, applying CycleDiffusion to text-to-image diffusion models, we show that large-scale text-to-image diffusion models can be used as zero-shot image-to-image editors. (2) One can guide pre-trained diffusion models and GANs by controlling the latent codes in a unified, plug-and-play formulation based on energy-based models. Using the CLIP model and a face recognition model as guidance, we demonstrate that diffusion models have better coverage of low-density sub-populations and individuals than GANs.*
|
||||
|
||||
*Tips*:
|
||||
- The Cycle Diffusion pipeline is fully compatible with any [Stable Diffusion](./stable_diffusion) checkpoints
|
||||
- Currently Cycle Diffusion only works with the [`DDIMScheduler`].
|
||||
|
||||
*Example*:
|
||||
|
||||
In the following we should how to best use the [`CycleDiffusionPipeline`]
|
||||
|
||||
```python
|
||||
import requests
|
||||
import torch
|
||||
from PIL import Image
|
||||
from io import BytesIO
|
||||
|
||||
from diffusers import CycleDiffusionPipeline, DDIMScheduler
|
||||
|
||||
# load the pipeline
|
||||
# make sure you're logged in with `huggingface-cli login`
|
||||
model_id_or_path = "CompVis/stable-diffusion-v1-4"
|
||||
scheduler = DDIMScheduler.from_pretrained(model_id_or_path, subfolder="scheduler")
|
||||
pipe = CycleDiffusionPipeline.from_pretrained(model_id_or_path, scheduler=scheduler).to("cuda")
|
||||
|
||||
# let's download an initial image
|
||||
url = "https://raw.githubusercontent.com/ChenWu98/cycle-diffusion/main/data/dalle2/An%20astronaut%20riding%20a%20horse.png"
|
||||
response = requests.get(url)
|
||||
init_image = Image.open(BytesIO(response.content)).convert("RGB")
|
||||
init_image = init_image.resize((512, 512))
|
||||
init_image.save("horse.png")
|
||||
|
||||
# let's specify a prompt
|
||||
source_prompt = "An astronaut riding a horse"
|
||||
prompt = "An astronaut riding an elephant"
|
||||
|
||||
# call the pipeline
|
||||
image = pipe(
|
||||
prompt=prompt,
|
||||
source_prompt=source_prompt,
|
||||
image=init_image,
|
||||
num_inference_steps=100,
|
||||
eta=0.1,
|
||||
strength=0.8,
|
||||
guidance_scale=2,
|
||||
source_guidance_scale=1,
|
||||
).images[0]
|
||||
|
||||
image.save("horse_to_elephant.png")
|
||||
|
||||
# let's try another example
|
||||
# See more samples at the original repo: https://github.com/ChenWu98/cycle-diffusion
|
||||
url = "https://raw.githubusercontent.com/ChenWu98/cycle-diffusion/main/data/dalle2/A%20black%20colored%20car.png"
|
||||
response = requests.get(url)
|
||||
init_image = Image.open(BytesIO(response.content)).convert("RGB")
|
||||
init_image = init_image.resize((512, 512))
|
||||
init_image.save("black.png")
|
||||
|
||||
source_prompt = "A black colored car"
|
||||
prompt = "A blue colored car"
|
||||
|
||||
# call the pipeline
|
||||
torch.manual_seed(0)
|
||||
image = pipe(
|
||||
prompt=prompt,
|
||||
source_prompt=source_prompt,
|
||||
image=init_image,
|
||||
num_inference_steps=100,
|
||||
eta=0.1,
|
||||
strength=0.85,
|
||||
guidance_scale=3,
|
||||
source_guidance_scale=1,
|
||||
).images[0]
|
||||
|
||||
image.save("black_to_blue.png")
|
||||
```
|
||||
|
||||
## CycleDiffusionPipeline
|
||||
[[autodoc]] CycleDiffusionPipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,59 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Scalable Diffusion Models with Transformers (DiT)
|
||||
|
||||
## Overview
|
||||
|
||||
[Scalable Diffusion Models with Transformers](https://arxiv.org/abs/2212.09748) (DiT) by William Peebles and Saining Xie.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*We explore a new class of diffusion models based on the transformer architecture. We train latent diffusion models of images, replacing the commonly-used U-Net backbone with a transformer that operates on latent patches. We analyze the scalability of our Diffusion Transformers (DiTs) through the lens of forward pass complexity as measured by Gflops. We find that DiTs with higher Gflops -- through increased transformer depth/width or increased number of input tokens -- consistently have lower FID. In addition to possessing good scalability properties, our largest DiT-XL/2 models outperform all prior diffusion models on the class-conditional ImageNet 512x512 and 256x256 benchmarks, achieving a state-of-the-art FID of 2.27 on the latter.*
|
||||
|
||||
The original codebase of this paper can be found here: [facebookresearch/dit](https://github.com/facebookresearch/dit).
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Colab
|
||||
|---|---|:---:|
|
||||
| [pipeline_dit.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/dit/pipeline_dit.py) | *Conditional Image Generation* | - |
|
||||
|
||||
|
||||
## Usage example
|
||||
|
||||
```python
|
||||
from diffusers import DiTPipeline, DPMSolverMultistepScheduler
|
||||
import torch
|
||||
|
||||
pipe = DiTPipeline.from_pretrained("facebook/DiT-XL-2-256", torch_dtype=torch.float16)
|
||||
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
# pick words from Imagenet class labels
|
||||
pipe.labels # to print all available words
|
||||
|
||||
# pick words that exist in ImageNet
|
||||
words = ["white shark", "umbrella"]
|
||||
|
||||
class_ids = pipe.get_label_ids(words)
|
||||
|
||||
generator = torch.manual_seed(33)
|
||||
output = pipe(class_labels=class_ids, num_inference_steps=25, generator=generator)
|
||||
|
||||
image = output.images[0] # label 'white shark'
|
||||
```
|
||||
|
||||
## DiTPipeline
|
||||
[[autodoc]] DiTPipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,74 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# PaintByExample
|
||||
|
||||
## Overview
|
||||
|
||||
[Paint by Example: Exemplar-based Image Editing with Diffusion Models](https://arxiv.org/abs/2211.13227) by Binxin Yang, Shuyang Gu, Bo Zhang, Ting Zhang, Xuejin Chen, Xiaoyan Sun, Dong Chen, Fang Wen
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*Language-guided image editing has achieved great success recently. In this paper, for the first time, we investigate exemplar-guided image editing for more precise control. We achieve this goal by leveraging self-supervised training to disentangle and re-organize the source image and the exemplar. However, the naive approach will cause obvious fusing artifacts. We carefully analyze it and propose an information bottleneck and strong augmentations to avoid the trivial solution of directly copying and pasting the exemplar image. Meanwhile, to ensure the controllability of the editing process, we design an arbitrary shape mask for the exemplar image and leverage the classifier-free guidance to increase the similarity to the exemplar image. The whole framework involves a single forward of the diffusion model without any iterative optimization. We demonstrate that our method achieves an impressive performance and enables controllable editing on in-the-wild images with high fidelity.*
|
||||
|
||||
The original codebase can be found [here](https://github.com/Fantasy-Studio/Paint-by-Example).
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Colab
|
||||
|---|---|:---:|
|
||||
| [pipeline_paint_by_example.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/paint_by_example/pipeline_paint_by_example.py) | *Image-Guided Image Painting* | - |
|
||||
|
||||
## Tips
|
||||
|
||||
- PaintByExample is supported by the official [Fantasy-Studio/Paint-by-Example](https://huggingface.co/Fantasy-Studio/Paint-by-Example) checkpoint. The checkpoint has been warm-started from the [CompVis/stable-diffusion-v1-4](https://huggingface.co/CompVis/stable-diffusion-v1-4) and with the objective to inpaint partly masked images conditioned on example / reference images
|
||||
- To quickly demo *PaintByExample*, please have a look at [this demo](https://huggingface.co/spaces/Fantasy-Studio/Paint-by-Example)
|
||||
- You can run the following code snippet as an example:
|
||||
|
||||
|
||||
```python
|
||||
# !pip install diffusers transformers
|
||||
|
||||
import PIL
|
||||
import requests
|
||||
import torch
|
||||
from io import BytesIO
|
||||
from diffusers import DiffusionPipeline
|
||||
|
||||
|
||||
def download_image(url):
|
||||
response = requests.get(url)
|
||||
return PIL.Image.open(BytesIO(response.content)).convert("RGB")
|
||||
|
||||
|
||||
img_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/image/example_1.png"
|
||||
mask_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/mask/example_1.png"
|
||||
example_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/reference/example_1.jpg"
|
||||
|
||||
init_image = download_image(img_url).resize((512, 512))
|
||||
mask_image = download_image(mask_url).resize((512, 512))
|
||||
example_image = download_image(example_url).resize((512, 512))
|
||||
|
||||
pipe = DiffusionPipeline.from_pretrained(
|
||||
"Fantasy-Studio/Paint-by-Example",
|
||||
torch_dtype=torch.float16,
|
||||
)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
image = pipe(image=init_image, mask_image=mask_image, example_image=example_image).images[0]
|
||||
image
|
||||
```
|
||||
|
||||
## PaintByExamplePipeline
|
||||
[[autodoc]] PaintByExamplePipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,77 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# RePaint
|
||||
|
||||
## Overview
|
||||
|
||||
[RePaint: Inpainting using Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2201.09865) (PNDM) by Andreas Lugmayr, Martin Danelljan, Andres Romero, Fisher Yu, Radu Timofte, Luc Van Gool.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
Free-form inpainting is the task of adding new content to an image in the regions specified by an arbitrary binary mask. Most existing approaches train for a certain distribution of masks, which limits their generalization capabilities to unseen mask types. Furthermore, training with pixel-wise and perceptual losses often leads to simple textural extensions towards the missing areas instead of semantically meaningful generation. In this work, we propose RePaint: A Denoising Diffusion Probabilistic Model (DDPM) based inpainting approach that is applicable to even extreme masks. We employ a pretrained unconditional DDPM as the generative prior. To condition the generation process, we only alter the reverse diffusion iterations by sampling the unmasked regions using the given image information. Since this technique does not modify or condition the original DDPM network itself, the model produces high-quality and diverse output images for any inpainting form. We validate our method for both faces and general-purpose image inpainting using standard and extreme masks.
|
||||
RePaint outperforms state-of-the-art Autoregressive, and GAN approaches for at least five out of six mask distributions.
|
||||
|
||||
The original codebase can be found [here](https://github.com/andreas128/RePaint).
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Colab
|
||||
|-------------------------------------------------------------------------------------------------------------------------------|--------------------|:---:|
|
||||
| [pipeline_repaint.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/repaint/pipeline_repaint.py) | *Image Inpainting* | - |
|
||||
|
||||
## Usage example
|
||||
|
||||
```python
|
||||
from io import BytesIO
|
||||
|
||||
import torch
|
||||
|
||||
import PIL
|
||||
import requests
|
||||
from diffusers import RePaintPipeline, RePaintScheduler
|
||||
|
||||
|
||||
def download_image(url):
|
||||
response = requests.get(url)
|
||||
return PIL.Image.open(BytesIO(response.content)).convert("RGB")
|
||||
|
||||
|
||||
img_url = "https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/repaint/celeba_hq_256.png"
|
||||
mask_url = "https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/repaint/mask_256.png"
|
||||
|
||||
# Load the original image and the mask as PIL images
|
||||
original_image = download_image(img_url).resize((256, 256))
|
||||
mask_image = download_image(mask_url).resize((256, 256))
|
||||
|
||||
# Load the RePaint scheduler and pipeline based on a pretrained DDPM model
|
||||
scheduler = RePaintScheduler.from_pretrained("google/ddpm-ema-celebahq-256")
|
||||
pipe = RePaintPipeline.from_pretrained("google/ddpm-ema-celebahq-256", scheduler=scheduler)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
generator = torch.Generator(device="cuda").manual_seed(0)
|
||||
output = pipe(
|
||||
original_image=original_image,
|
||||
mask_image=mask_image,
|
||||
num_inference_steps=250,
|
||||
eta=0.0,
|
||||
jump_length=10,
|
||||
jump_n_sample=10,
|
||||
generator=generator,
|
||||
)
|
||||
inpainted_image = output.images[0]
|
||||
```
|
||||
|
||||
## RePaintPipeline
|
||||
[[autodoc]] RePaintPipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,79 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Semantic Guidance
|
||||
|
||||
Semantic Guidance for Diffusion Models was proposed in [SEGA: Instructing Diffusion using Semantic Dimensions](https://arxiv.org/abs/2301.12247) and provides strong semantic control over the image generation.
|
||||
Small changes to the text prompt usually result in entirely different output images. However, with SEGA a variety of changes to the image are enabled that can be controlled easily and intuitively, and stay true to the original image composition.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*Text-to-image diffusion models have recently received a lot of interest for their astonishing ability to produce high-fidelity images from text only. However, achieving one-shot generation that aligns with the user's intent is nearly impossible, yet small changes to the input prompt often result in very different images. This leaves the user with little semantic control. To put the user in control, we show how to interact with the diffusion process to flexibly steer it along semantic directions. This semantic guidance (SEGA) allows for subtle and extensive edits, changes in composition and style, as well as optimizing the overall artistic conception. We demonstrate SEGA's effectiveness on a variety of tasks and provide evidence for its versatility and flexibility.*
|
||||
|
||||
|
||||
*Overview*:
|
||||
|
||||
| Pipeline | Tasks | Colab | Demo
|
||||
|---|---|:---:|:---:|
|
||||
| [pipeline_semantic_stable_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/semantic_stable_diffusion/pipeline_semantic_stable_diffusion) | *Text-to-Image Generation* | [](https://colab.research.google.com/github/ml-research/semantic-image-editing/blob/main/examples/SemanticGuidance.ipynb) | [Coming Soon](https://huggingface.co/AIML-TUDA)
|
||||
|
||||
## Tips
|
||||
|
||||
- The Semantic Guidance pipeline can be used with any [Stable Diffusion](./api/pipelines/stable_diffusion/text2img) checkpoint.
|
||||
|
||||
### Run Semantic Guidance
|
||||
|
||||
The interface of [`SemanticStableDiffusionPipeline`] provides several additional parameters to influence the image generation.
|
||||
Exemplary usage may look like this:
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import SemanticStableDiffusionPipeline
|
||||
|
||||
pipe = SemanticStableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
out = pipe(
|
||||
prompt="a photo of the face of a woman",
|
||||
num_images_per_prompt=1,
|
||||
guidance_scale=7,
|
||||
editing_prompt=[
|
||||
"smiling, smile", # Concepts to apply
|
||||
"glasses, wearing glasses",
|
||||
"curls, wavy hair, curly hair",
|
||||
"beard, full beard, mustache",
|
||||
],
|
||||
reverse_editing_direction=[False, False, False, False], # Direction of guidance i.e. increase all concepts
|
||||
edit_warmup_steps=[10, 10, 10, 10], # Warmup period for each concept
|
||||
edit_guidance_scale=[4, 5, 5, 5.4], # Guidance scale for each concept
|
||||
edit_threshold=[
|
||||
0.99,
|
||||
0.975,
|
||||
0.925,
|
||||
0.96,
|
||||
], # Threshold for each concept. Threshold equals the percentile of the latent space that will be discarded. I.e. threshold=0.99 uses 1% of the latent dimensions
|
||||
edit_momentum_scale=0.3, # Momentum scale that will be added to the latent guidance
|
||||
edit_mom_beta=0.6, # Momentum beta
|
||||
edit_weights=[1, 1, 1, 1, 1], # Weights of the individual concepts against each other
|
||||
)
|
||||
```
|
||||
|
||||
For more examples check the colab notebook.
|
||||
|
||||
## StableDiffusionSafePipelineOutput
|
||||
[[autodoc]] pipelines.semantic_stable_diffusion.SemanticStableDiffusionPipelineOutput
|
||||
- all
|
||||
|
||||
## SemanticStableDiffusionPipeline
|
||||
[[autodoc]] SemanticStableDiffusionPipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,75 +0,0 @@
|
||||
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Attend and Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models
|
||||
|
||||
## Overview
|
||||
|
||||
Attend and Excite for Stable Diffusion was proposed in [Attend-and-Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models](https://attendandexcite.github.io/Attend-and-Excite/) and provides textual attention control over the image generation.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*Text-to-image diffusion models have recently received a lot of interest for their astonishing ability to produce high-fidelity images from text only. However, achieving one-shot generation that aligns with the user's intent is nearly impossible, yet small changes to the input prompt often result in very different images. This leaves the user with little semantic control. To put the user in control, we show how to interact with the diffusion process to flexibly steer it along semantic directions. This semantic guidance (SEGA) allows for subtle and extensive edits, changes in composition and style, as well as optimizing the overall artistic conception. We demonstrate SEGA's effectiveness on a variety of tasks and provide evidence for its versatility and flexibility.*
|
||||
|
||||
Resources
|
||||
|
||||
* [Project Page](https://attendandexcite.github.io/Attend-and-Excite/)
|
||||
* [Paper](https://arxiv.org/abs/2301.13826)
|
||||
* [Original Code](https://github.com/AttendAndExcite/Attend-and-Excite)
|
||||
* [Demo](https://huggingface.co/spaces/AttendAndExcite/Attend-and-Excite)
|
||||
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Colab | Demo
|
||||
|---|---|:---:|:---:|
|
||||
| [pipeline_semantic_stable_diffusion_attend_and_excite.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_semantic_stable_diffusion_attend_and_excite) | *Text-to-Image Generation* | - | -
|
||||
|
||||
|
||||
### Usage example
|
||||
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import StableDiffusionAttendAndExcitePipeline
|
||||
|
||||
model_id = "CompVis/stable-diffusion-v1-4"
|
||||
pipe = StableDiffusionAttendAndExcitePipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "a cat and a frog"
|
||||
|
||||
# use get_indices function to find out indices of the tokens you want to alter
|
||||
pipe.get_indices(prompt)
|
||||
|
||||
token_indices = [2, 5]
|
||||
seed = 6141
|
||||
generator = torch.Generator("cuda").manual_seed(seed)
|
||||
|
||||
images = pipe(
|
||||
prompt=prompt,
|
||||
token_indices=token_indices,
|
||||
guidance_scale=7.5,
|
||||
generator=generator,
|
||||
num_inference_steps=50,
|
||||
max_iter_to_alter=25,
|
||||
).images
|
||||
|
||||
image = images[0]
|
||||
image.save(f"../images/{prompt}_{seed}.png")
|
||||
```
|
||||
|
||||
|
||||
## StableDiffusionAttendAndExcitePipeline
|
||||
[[autodoc]] StableDiffusionAttendAndExcitePipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,33 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Depth-to-Image Generation
|
||||
|
||||
## StableDiffusionDepth2ImgPipeline
|
||||
|
||||
The depth-guided stable diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), and [LAION](https://laion.ai/), as part of Stable Diffusion 2.0. It uses [MiDas](https://github.com/isl-org/MiDaS) to infer depth based on an image.
|
||||
|
||||
[`StableDiffusionDepth2ImgPipeline`] lets you pass a text prompt and an initial image to condition the generation of new images as well as a `depth_map` to preserve the images’ structure.
|
||||
|
||||
The original codebase can be found here:
|
||||
- *Stable Diffusion v2*: [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion#depth-conditional-stable-diffusion)
|
||||
|
||||
Available Checkpoints are:
|
||||
- *stable-diffusion-2-depth*: [stabilityai/stable-diffusion-2-depth](https://huggingface.co/stabilityai/stable-diffusion-2-depth)
|
||||
|
||||
[[autodoc]] StableDiffusionDepth2ImgPipeline
|
||||
- all
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
@@ -1,31 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Image Variation
|
||||
|
||||
## StableDiffusionImageVariationPipeline
|
||||
|
||||
[`StableDiffusionImageVariationPipeline`] lets you generate variations from an input image using Stable Diffusion. It uses a fine-tuned version of Stable Diffusion model, trained by [Justin Pinkney](https://www.justinpinkney.com/) (@Buntworthy) at [Lambda](https://lambdalabs.com/)
|
||||
|
||||
The original codebase can be found here:
|
||||
[Stable Diffusion Image Variations](https://github.com/LambdaLabsML/lambda-diffusers#stable-diffusion-image-variations)
|
||||
|
||||
Available Checkpoints are:
|
||||
- *sd-image-variations-diffusers*: [lambdalabs/sd-image-variations-diffusers](https://huggingface.co/lambdalabs/sd-image-variations-diffusers)
|
||||
|
||||
[[autodoc]] StableDiffusionImageVariationPipeline
|
||||
- all
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
@@ -1,29 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Image-to-Image Generation
|
||||
|
||||
## StableDiffusionImg2ImgPipeline
|
||||
|
||||
The Stable Diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionImg2ImgPipeline`] lets you pass a text prompt and an initial image to condition the generation of new images using Stable Diffusion.
|
||||
|
||||
The original codebase can be found here: [CampVis/stable-diffusion](https://github.com/CompVis/stable-diffusion/blob/main/scripts/img2img.py)
|
||||
|
||||
[`StableDiffusionImg2ImgPipeline`] is compatible with all Stable Diffusion checkpoints for [Text-to-Image](./text2img)
|
||||
|
||||
[[autodoc]] StableDiffusionImg2ImgPipeline
|
||||
- all
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
@@ -1,33 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Text-Guided Image Inpainting
|
||||
|
||||
## StableDiffusionInpaintPipeline
|
||||
|
||||
The Stable Diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionInpaintPipeline`] lets you edit specific parts of an image by providing a mask and a text prompt using Stable Diffusion.
|
||||
|
||||
The original codebase can be found here:
|
||||
- *Stable Diffusion V1*: [CampVis/stable-diffusion](https://github.com/runwayml/stable-diffusion#inpainting-with-stable-diffusion)
|
||||
- *Stable Diffusion V2*: [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion#image-inpainting-with-stable-diffusion)
|
||||
|
||||
Available checkpoints are:
|
||||
- *stable-diffusion-inpainting (512x512 resolution)*: [runwayml/stable-diffusion-inpainting](https://huggingface.co/runwayml/stable-diffusion-inpainting)
|
||||
- *stable-diffusion-2-inpainting (512x512 resolution)*: [stabilityai/stable-diffusion-2-inpainting](https://huggingface.co/stabilityai/stable-diffusion-2-inpainting)
|
||||
|
||||
[[autodoc]] StableDiffusionInpaintPipeline
|
||||
- all
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
@@ -1,33 +0,0 @@
|
||||
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Stable Diffusion Latent Upscaler
|
||||
|
||||
## StableDiffusionLatentUpscalePipeline
|
||||
|
||||
The Stable Diffusion Latent Upscaler model was created by [Katherine Crowson](https://github.com/crowsonkb/k-diffusion) in collaboration with [Stability AI](https://stability.ai/). It can be used on top of any [`StableDiffusionUpscalePipeline`] checkpoint to enhance its output image resolution by a factor of 2.
|
||||
|
||||
A notebook that demonstrates the original implementation can be found here:
|
||||
- [Stable Diffusion Upscaler Demo](https://colab.research.google.com/drive/1o1qYJcFeywzCIdkfKJy7cTpgZTCM2EI4)
|
||||
|
||||
Available Checkpoints are:
|
||||
- *stabilityai/latent-upscaler*: [stabilityai/sd-x2-latent-upscaler](https://huggingface.co/stabilityai/sd-x2-latent-upscaler)
|
||||
|
||||
|
||||
[[autodoc]] StableDiffusionLatentUpscalePipeline
|
||||
- all
|
||||
- __call__
|
||||
- enable_sequential_cpu_offload
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
@@ -1,81 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Stable diffusion pipelines
|
||||
|
||||
Stable Diffusion is a text-to-image _latent diffusion_ model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/). It's trained on 512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/) dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs.
|
||||
|
||||
Latent diffusion is the research on top of which Stable Diffusion was built. It was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://arxiv.org/abs/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer. You can learn more details about it in the [specific pipeline for latent diffusion](pipelines/latent_diffusion) that is part of 🤗 Diffusers.
|
||||
|
||||
For more details about how Stable Diffusion works and how it differs from the base latent diffusion model, please refer to the official [launch announcement post](https://stability.ai/blog/stable-diffusion-announcement) and [this section of our own blog post](https://huggingface.co/blog/stable_diffusion#how-does-stable-diffusion-work).
|
||||
|
||||
*Tips*:
|
||||
- To tweak your prompts on a specific result you liked, you can generate your own latents, as demonstrated in the following notebook: [](https://colab.research.google.com/github/pcuenca/diffusers-examples/blob/main/notebooks/stable-diffusion-seeds.ipynb)
|
||||
|
||||
*Overview*:
|
||||
|
||||
| Pipeline | Tasks | Colab | Demo
|
||||
|---|---|:---:|:---:|
|
||||
| [StableDiffusionPipeline](./text2img) | *Text-to-Image Generation* | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb) | [🤗 Stable Diffusion](https://huggingface.co/spaces/stabilityai/stable-diffusion)
|
||||
| [StableDiffusionImg2ImgPipeline](./img2img) | *Image-to-Image Text-Guided Generation* | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb) | [🤗 Diffuse the Rest](https://huggingface.co/spaces/huggingface/diffuse-the-rest)
|
||||
| [StableDiffusionInpaintPipeline](./inpaint) | **Experimental** – *Text-Guided Image Inpainting* | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb) | Coming soon
|
||||
| [StableDiffusionDepth2ImgPipeline](./depth2img) | **Experimental** – *Depth-to-Image Text-Guided Generation * | | Coming soon
|
||||
| [StableDiffusionImageVariationPipeline](./image_variation) | **Experimental** – *Image Variation Generation * | | [🤗 Stable Diffusion Image Variations](https://huggingface.co/spaces/lambdalabs/stable-diffusion-image-variations)
|
||||
| [StableDiffusionUpscalePipeline](./upscale) | **Experimental** – *Text-Guided Image Super-Resolution * | | Coming soon
|
||||
| [StableDiffusionLatentUpscalePipeline](./latent_upscale) | **Experimental** – *Text-Guided Image Super-Resolution * | | Coming soon
|
||||
| [StableDiffusionInstructPix2PixPipeline](./pix2pix) | **Experimental** – *Text-Based Image Editing * | | [InstructPix2Pix: Learning to Follow Image Editing Instructions](https://huggingface.co/spaces/timbrooks/instruct-pix2pix)
|
||||
| [StableDiffusionAttendAndExcitePipeline](./attend_and_excite) | **Experimental** – *Text-to-Image Generation * | | [Attend-and-Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models](https://huggingface.co/spaces/AttendAndExcite/Attend-and-Excite)
|
||||
| [StableDiffusionPix2PixZeroPipeline](./pix2pix_zero) | **Experimental** – *Text-Based Image Editing * | | [Zero-shot Image-to-Image Translation](https://arxiv.org/abs/2302.03027)
|
||||
|
||||
|
||||
|
||||
## Tips
|
||||
|
||||
### How to load and use different schedulers.
|
||||
|
||||
The stable diffusion pipeline uses [`PNDMScheduler`] scheduler by default. But `diffusers` provides many other schedulers that can be used with the stable diffusion pipeline such as [`DDIMScheduler`], [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`] etc.
|
||||
To use a different scheduler, you can either change it via the [`ConfigMixin.from_config`] method or pass the `scheduler` argument to the `from_pretrained` method of the pipeline. For example, to use the [`EulerDiscreteScheduler`], you can do the following:
|
||||
|
||||
```python
|
||||
>>> from diffusers import StableDiffusionPipeline, EulerDiscreteScheduler
|
||||
|
||||
>>> pipeline = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4")
|
||||
>>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
|
||||
|
||||
>>> # or
|
||||
>>> euler_scheduler = EulerDiscreteScheduler.from_pretrained("CompVis/stable-diffusion-v1-4", subfolder="scheduler")
|
||||
>>> pipeline = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", scheduler=euler_scheduler)
|
||||
```
|
||||
|
||||
|
||||
### How to convert all use cases with multiple or single pipeline
|
||||
|
||||
If you want to use all possible use cases in a single `DiffusionPipeline` you can either:
|
||||
- Make use of the [Stable Diffusion Mega Pipeline](https://github.com/huggingface/diffusers/tree/main/examples/community#stable-diffusion-mega) or
|
||||
- Make use of the `components` functionality to instantiate all components in the most memory-efficient way:
|
||||
|
||||
```python
|
||||
>>> from diffusers import (
|
||||
... StableDiffusionPipeline,
|
||||
... StableDiffusionImg2ImgPipeline,
|
||||
... StableDiffusionInpaintPipeline,
|
||||
... )
|
||||
|
||||
>>> text2img = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4")
|
||||
>>> img2img = StableDiffusionImg2ImgPipeline(**text2img.components)
|
||||
>>> inpaint = StableDiffusionInpaintPipeline(**text2img.components)
|
||||
|
||||
>>> # now you can use text2img(...), img2img(...), inpaint(...) just like the call methods of each respective pipeline
|
||||
```
|
||||
|
||||
## StableDiffusionPipelineOutput
|
||||
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
|
||||
@@ -1,58 +0,0 @@
|
||||
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation
|
||||
|
||||
## Overview
|
||||
|
||||
[MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation](https://arxiv.org/abs/2302.08113) by Omer Bar-Tal, Lior Yariv, Yaron Lipman, and Tali Dekel.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*Recent advances in text-to-image generation with diffusion models present transformative capabilities in image quality. However, user controllability of the generated image, and fast adaptation to new tasks still remains an open challenge, currently mostly addressed by costly and long re-training and fine-tuning or ad-hoc adaptations to specific image generation tasks. In this work, we present MultiDiffusion, a unified framework that enables versatile and controllable image generation, using a pre-trained text-to-image diffusion model, without any further training or finetuning. At the center of our approach is a new generation process, based on an optimization task that binds together multiple diffusion generation processes with a shared set of parameters or constraints. We show that MultiDiffusion can be readily applied to generate high quality and diverse images that adhere to user-provided controls, such as desired aspect ratio (e.g., panorama), and spatial guiding signals, ranging from tight segmentation masks to bounding boxes.
|
||||
|
||||
Resources:
|
||||
|
||||
* [Project Page](https://multidiffusion.github.io/).
|
||||
* [Paper](https://arxiv.org/abs/2302.08113).
|
||||
* [Original Code](https://github.com/omerbt/MultiDiffusion).
|
||||
* [Demo](https://huggingface.co/spaces/weizmannscience/MultiDiffusion).
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Demo
|
||||
|---|---|:---:|
|
||||
| [StableDiffusionPanoramaPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_panorama.py) | *Text-Guided Panorama View Generation* | [🤗 Space](https://huggingface.co/spaces/weizmannscience/MultiDiffusion)) |
|
||||
|
||||
<!-- TODO: add Colab -->
|
||||
|
||||
## Usage example
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import StableDiffusionPanoramaPipeline, DDIMScheduler
|
||||
|
||||
model_ckpt = "stabilityai/stable-diffusion-2-base"
|
||||
scheduler = DDIMScheduler.from_pretrained(model_ckpt, subfolder="scheduler")
|
||||
pipe = StableDiffusionPanoramaPipeline.from_pretrained(model_ckpt, scheduler=scheduler, torch_dtype=torch.float16)
|
||||
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "a photo of the dolomites"
|
||||
image = pipe(prompt).images[0]
|
||||
image.save("dolomites.png")
|
||||
```
|
||||
|
||||
## StableDiffusionPanoramaPipeline
|
||||
[[autodoc]] StableDiffusionPanoramaPipeline
|
||||
- __call__
|
||||
- all
|
||||
@@ -1,70 +0,0 @@
|
||||
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# InstructPix2Pix: Learning to Follow Image Editing Instructions
|
||||
|
||||
## Overview
|
||||
|
||||
[InstructPix2Pix: Learning to Follow Image Editing Instructions](https://arxiv.org/abs/2211.09800) by Tim Brooks, Aleksander Holynski and Alexei A. Efros.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*We propose a method for editing images from human instructions: given an input image and a written instruction that tells the model what to do, our model follows these instructions to edit the image. To obtain training data for this problem, we combine the knowledge of two large pretrained models -- a language model (GPT-3) and a text-to-image model (Stable Diffusion) -- to generate a large dataset of image editing examples. Our conditional diffusion model, InstructPix2Pix, is trained on our generated data, and generalizes to real images and user-written instructions at inference time. Since it performs edits in the forward pass and does not require per example fine-tuning or inversion, our model edits images quickly, in a matter of seconds. We show compelling editing results for a diverse collection of input images and written instructions.*
|
||||
|
||||
Resources:
|
||||
|
||||
* [Project Page](https://www.timothybrooks.com/instruct-pix2pix).
|
||||
* [Paper](https://arxiv.org/abs/2211.09800).
|
||||
* [Original Code](https://github.com/timothybrooks/instruct-pix2pix).
|
||||
* [Demo](https://huggingface.co/spaces/timbrooks/instruct-pix2pix).
|
||||
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Demo
|
||||
|---|---|:---:|
|
||||
| [StableDiffusionInstructPix2PixPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_instruct_pix2pix.py) | *Text-Based Image Editing* | [🤗 Space](https://huggingface.co/spaces/timbrooks/instruct-pix2pix) |
|
||||
|
||||
<!-- TODO: add Colab -->
|
||||
|
||||
## Usage example
|
||||
|
||||
```python
|
||||
import PIL
|
||||
import requests
|
||||
import torch
|
||||
from diffusers import StableDiffusionInstructPix2PixPipeline
|
||||
|
||||
model_id = "timbrooks/instruct-pix2pix"
|
||||
pipe = StableDiffusionInstructPix2PixPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
|
||||
|
||||
url = "https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/mountain.png"
|
||||
|
||||
|
||||
def download_image(url):
|
||||
image = PIL.Image.open(requests.get(url, stream=True).raw)
|
||||
image = PIL.ImageOps.exif_transpose(image)
|
||||
image = image.convert("RGB")
|
||||
return image
|
||||
|
||||
|
||||
image = download_image(url)
|
||||
|
||||
prompt = "make the mountains snowy"
|
||||
edit = pipe(prompt, image=image, num_inference_steps=20, image_guidance_scale=1.5, guidance_scale=7).images[0]
|
||||
images[0].save("snowy_mountains.png")
|
||||
```
|
||||
|
||||
## StableDiffusionInstructPix2PixPipeline
|
||||
[[autodoc]] StableDiffusionInstructPix2PixPipeline
|
||||
- __call__
|
||||
- all
|
||||
@@ -1,289 +0,0 @@
|
||||
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Zero-shot Image-to-Image Translation
|
||||
|
||||
## Overview
|
||||
|
||||
[Zero-shot Image-to-Image Translation](https://arxiv.org/abs/2302.03027).
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*Large-scale text-to-image generative models have shown their remarkable ability to synthesize diverse and high-quality images. However, it is still challenging to directly apply these models for editing real images for two reasons. First, it is hard for users to come up with a perfect text prompt that accurately describes every visual detail in the input image. Second, while existing models can introduce desirable changes in certain regions, they often dramatically alter the input content and introduce unexpected changes in unwanted regions. In this work, we propose pix2pix-zero, an image-to-image translation method that can preserve the content of the original image without manual prompting. We first automatically discover editing directions that reflect desired edits in the text embedding space. To preserve the general content structure after editing, we further propose cross-attention guidance, which aims to retain the cross-attention maps of the input image throughout the diffusion process. In addition, our method does not need additional training for these edits and can directly use the existing pre-trained text-to-image diffusion model. We conduct extensive experiments and show that our method outperforms existing and concurrent works for both real and synthetic image editing.*
|
||||
|
||||
Resources:
|
||||
|
||||
* [Project Page](https://pix2pixzero.github.io/).
|
||||
* [Paper](https://arxiv.org/abs/2302.03027).
|
||||
* [Original Code](https://github.com/pix2pixzero/pix2pix-zero).
|
||||
|
||||
## Tips
|
||||
|
||||
* The pipeline can be conditioned on real input images. Check out the code examples below to know more.
|
||||
* The pipeline exposes two arguments namely `source_embeds` and `target_embeds`
|
||||
that let you control the direction of the semantic edits in the final image to be generated. Let's say,
|
||||
you wanted to translate from "cat" to "dog". In this case, the edit direction will be "cat -> dog". To reflect
|
||||
this in the pipeline, you simply have to set the embeddings related to the phrases including "cat" to
|
||||
`source_embeds` and "dog" to `target_embeds`. Refer to the code example below for more details.
|
||||
* When you're using this pipeline from a prompt, specify the _source_ concept in the prompt. Taking
|
||||
the above example, a valid input prompt would be: "a high resolution painting of a **cat** in the style of van gough".
|
||||
* If you wanted to reverse the direction in the example above, i.e., "dog -> cat", then it's recommended to:
|
||||
* Swap the `source_embeds` and `target_embeds`.
|
||||
* Change the input prompt to include "dog".
|
||||
* To learn more about how the source and target embeddings are generated, refer to the [original
|
||||
paper](https://arxiv.org/abs/2302.03027). Below, we also provide some directions on how to generate the embeddings.
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Demo
|
||||
|---|---|:---:|
|
||||
| [StableDiffusionPix2PixZeroPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_pix2pix_zero.py) | *Text-Based Image Editing* | [🤗 Space] (soon) |
|
||||
|
||||
<!-- TODO: add Colab -->
|
||||
|
||||
## Usage example
|
||||
|
||||
### Based on an image generated with the input prompt
|
||||
|
||||
```python
|
||||
import requests
|
||||
import torch
|
||||
|
||||
from diffusers import DDIMScheduler, StableDiffusionPix2PixZeroPipeline
|
||||
|
||||
|
||||
def download(embedding_url, local_filepath):
|
||||
r = requests.get(embedding_url)
|
||||
with open(local_filepath, "wb") as f:
|
||||
f.write(r.content)
|
||||
|
||||
|
||||
model_ckpt = "CompVis/stable-diffusion-v1-4"
|
||||
pipeline = StableDiffusionPix2PixZeroPipeline.from_pretrained(
|
||||
model_ckpt, conditions_input_image=False, torch_dtype=torch.float16
|
||||
)
|
||||
pipeline.scheduler = DDIMScheduler.from_config(pipeline.scheduler.config)
|
||||
pipeline.to("cuda")
|
||||
|
||||
prompt = "a high resolution painting of a cat in the style of van gough"
|
||||
src_embs_url = "https://github.com/pix2pixzero/pix2pix-zero/raw/main/assets/embeddings_sd_1.4/cat.pt"
|
||||
target_embs_url = "https://github.com/pix2pixzero/pix2pix-zero/raw/main/assets/embeddings_sd_1.4/dog.pt"
|
||||
|
||||
for url in [src_embs_url, target_embs_url]:
|
||||
download(url, url.split("/")[-1])
|
||||
|
||||
src_embeds = torch.load(src_embs_url.split("/")[-1])
|
||||
target_embeds = torch.load(target_embs_url.split("/")[-1])
|
||||
|
||||
images = pipeline(
|
||||
prompt,
|
||||
source_embeds=src_embeds,
|
||||
target_embeds=target_embeds,
|
||||
num_inference_steps=50,
|
||||
cross_attention_guidance_amount=0.15,
|
||||
).images
|
||||
images[0].save("edited_image_dog.png")
|
||||
```
|
||||
|
||||
### Based on an input image
|
||||
|
||||
When the pipeline is conditioned on an input image, we first obtain an inverted
|
||||
noise from it using a `DDIMInverseScheduler` with the help of a generated caption. Then
|
||||
the inverted noise is used to start the generation process.
|
||||
|
||||
First, let's load our pipeline:
|
||||
|
||||
```py
|
||||
import torch
|
||||
from transformers import BlipForConditionalGeneration, BlipProcessor
|
||||
from diffusers import DDIMScheduler, DDIMInverseScheduler, StableDiffusionPix2PixZeroPipeline
|
||||
|
||||
captioner_id = "Salesforce/blip-image-captioning-base"
|
||||
processor = BlipProcessor.from_pretrained(captioner_id)
|
||||
model = BlipForConditionalGeneration.from_pretrained(captioner_id, torch_dtype=torch.float16, low_cpu_mem_usage=True)
|
||||
|
||||
sd_model_ckpt = "CompVis/stable-diffusion-v1-4"
|
||||
pipeline = StableDiffusionPix2PixZeroPipeline.from_pretrained(
|
||||
sd_model_ckpt,
|
||||
caption_generator=model,
|
||||
caption_processor=processor,
|
||||
torch_dtype=torch.float16,
|
||||
safety_checker=None,
|
||||
)
|
||||
pipeline.scheduler = DDIMScheduler.from_config(pipeline.scheduler.config)
|
||||
pipeline.inverse_scheduler = DDIMInverseScheduler.from_config(pipeline.scheduler.config)
|
||||
pipeline.enable_model_cpu_offload()
|
||||
```
|
||||
|
||||
Then, we load an input image for conditioning and obtain a suitable caption for it:
|
||||
|
||||
```py
|
||||
import requests
|
||||
from PIL import Image
|
||||
|
||||
img_url = "https://github.com/pix2pixzero/pix2pix-zero/raw/main/assets/test_images/cats/cat_6.png"
|
||||
raw_image = Image.open(requests.get(img_url, stream=True).raw).convert("RGB").resize((512, 512))
|
||||
caption = pipeline.generate_caption(raw_image)
|
||||
```
|
||||
|
||||
Then we employ the generated caption and the input image to get the inverted noise:
|
||||
|
||||
```py
|
||||
generator = torch.manual_seed(0)
|
||||
inv_latents = pipeline.invert(caption, image=raw_image, generator=generator).latents
|
||||
```
|
||||
|
||||
Now, generate the image with edit directions:
|
||||
|
||||
```py
|
||||
# See the "Generating source and target embeddings" section below to
|
||||
# automate the generation of these captions with a pre-trained model like Flan-T5 as explained below.
|
||||
source_prompts = ["a cat sitting on the street", "a cat playing in the field", "a face of a cat"]
|
||||
target_prompts = ["a dog sitting on the street", "a dog playing in the field", "a face of a dog"]
|
||||
|
||||
source_embeds = pipeline.get_embeds(source_prompts, batch_size=2)
|
||||
target_embeds = pipeline.get_embeds(target_prompts, batch_size=2)
|
||||
|
||||
|
||||
image = pipeline(
|
||||
caption,
|
||||
source_embeds=source_embeds,
|
||||
target_embeds=target_embeds,
|
||||
num_inference_steps=50,
|
||||
cross_attention_guidance_amount=0.15,
|
||||
generator=generator,
|
||||
latents=inv_latents,
|
||||
negative_prompt=caption,
|
||||
).images[0]
|
||||
image.save("edited_image.png")
|
||||
```
|
||||
|
||||
## Generating source and target embeddings
|
||||
|
||||
The authors originally used the [GPT-3 API](https://openai.com/api/) to generate the source and target captions for discovering
|
||||
edit directions. However, we can also leverage open source and public models for the same purpose.
|
||||
Below, we provide an end-to-end example with the [Flan-T5](https://huggingface.co/docs/transformers/model_doc/flan-t5) model
|
||||
for generating captions and [CLIP](https://huggingface.co/docs/transformers/model_doc/clip) for
|
||||
computing embeddings on the generated captions.
|
||||
|
||||
**1. Load the generation model**:
|
||||
|
||||
```py
|
||||
import torch
|
||||
from transformers import AutoTokenizer, T5ForConditionalGeneration
|
||||
|
||||
tokenizer = AutoTokenizer.from_pretrained("google/flan-t5-xl")
|
||||
model = T5ForConditionalGeneration.from_pretrained("google/flan-t5-xl", device_map="auto", torch_dtype=torch.float16)
|
||||
```
|
||||
|
||||
**2. Construct a starting prompt**:
|
||||
|
||||
```py
|
||||
source_concept = "cat"
|
||||
target_concept = "dog"
|
||||
|
||||
source_text = f"Provide a caption for images containing a {source_concept}. "
|
||||
"The captions should be in English and should be no longer than 150 characters."
|
||||
|
||||
target_text = f"Provide a caption for images containing a {target_concept}. "
|
||||
"The captions should be in English and should be no longer than 150 characters."
|
||||
```
|
||||
|
||||
Here, we're interested in the "cat -> dog" direction.
|
||||
|
||||
**3. Generate captions**:
|
||||
|
||||
We can use a utility like so for this purpose.
|
||||
|
||||
```py
|
||||
def generate_captions(input_prompt):
|
||||
input_ids = tokenizer(input_prompt, return_tensors="pt").input_ids.to("cuda")
|
||||
|
||||
outputs = model.generate(
|
||||
input_ids, temperature=0.8, num_return_sequences=16, do_sample=True, max_new_tokens=128, top_k=10
|
||||
)
|
||||
return tokenizer.batch_decode(outputs, skip_special_tokens=True)
|
||||
```
|
||||
|
||||
And then we just call it to generate our captions:
|
||||
|
||||
```py
|
||||
source_captions = generate_captions(source_text)
|
||||
target_captions = generate_captions(target_concept)
|
||||
```
|
||||
|
||||
We encourage you to play around with the different parameters supported by the
|
||||
`generate()` method ([documentation](https://huggingface.co/docs/transformers/main/en/main_classes/text_generation#transformers.generation_tf_utils.TFGenerationMixin.generate)) for the generation quality you are looking for.
|
||||
|
||||
**4. Load the embedding model**:
|
||||
|
||||
Here, we need to use the same text encoder model used by the subsequent Stable Diffusion model.
|
||||
|
||||
```py
|
||||
from diffusers import StableDiffusionPix2PixZeroPipeline
|
||||
|
||||
pipeline = StableDiffusionPix2PixZeroPipeline.from_pretrained(
|
||||
"CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16
|
||||
)
|
||||
pipeline = pipeline.to("cuda")
|
||||
tokenizer = pipeline.tokenizer
|
||||
text_encoder = pipeline.text_encoder
|
||||
```
|
||||
|
||||
**5. Compute embeddings**:
|
||||
|
||||
```py
|
||||
import torch
|
||||
|
||||
def embed_captions(sentences, tokenizer, text_encoder, device="cuda"):
|
||||
with torch.no_grad():
|
||||
embeddings = []
|
||||
for sent in sentences:
|
||||
text_inputs = tokenizer(
|
||||
sent,
|
||||
padding="max_length",
|
||||
max_length=tokenizer.model_max_length,
|
||||
truncation=True,
|
||||
return_tensors="pt",
|
||||
)
|
||||
text_input_ids = text_inputs.input_ids
|
||||
prompt_embeds = text_encoder(text_input_ids.to(device), attention_mask=None)[0]
|
||||
embeddings.append(prompt_embeds)
|
||||
return torch.concatenate(embeddings, dim=0).mean(dim=0).unsqueeze(0)
|
||||
|
||||
source_embeddings = embed_captions(source_captions, tokenizer, text_encoder)
|
||||
target_embeddings = embed_captions(target_captions, tokenizer, text_encoder)
|
||||
```
|
||||
|
||||
And you're done! [Here](https://colab.research.google.com/drive/1tz2C1EdfZYAPlzXXbTnf-5PRBiR8_R1F?usp=sharing) is a Colab Notebook that you can use to interact with the entire process.
|
||||
|
||||
Now, you can use these embeddings directly while calling the pipeline:
|
||||
|
||||
```py
|
||||
from diffusers import DDIMScheduler
|
||||
|
||||
pipeline.scheduler = DDIMScheduler.from_config(pipeline.scheduler.config)
|
||||
|
||||
images = pipeline(
|
||||
prompt,
|
||||
source_embeds=source_embeddings,
|
||||
target_embeds=target_embeddings,
|
||||
num_inference_steps=50,
|
||||
cross_attention_guidance_amount=0.15,
|
||||
).images
|
||||
images[0].save("edited_image_dog.png")
|
||||
```
|
||||
|
||||
## StableDiffusionPix2PixZeroPipeline
|
||||
[[autodoc]] StableDiffusionPix2PixZeroPipeline
|
||||
- __call__
|
||||
- all
|
||||
@@ -1,64 +0,0 @@
|
||||
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Self-Attention Guidance (SAG)
|
||||
|
||||
## Overview
|
||||
|
||||
[Self-Attention Guidance](https://arxiv.org/abs/2210.00939) by Susung Hong et al.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*Denoising diffusion models (DDMs) have been drawing much attention for their appreciable sample quality and diversity. Despite their remarkable performance, DDMs remain black boxes on which further study is necessary to take a profound step. Motivated by this, we delve into the design of conventional U-shaped diffusion models. More specifically, we investigate the self-attention modules within these models through carefully designed experiments and explore their characteristics. In addition, inspired by the studies that substantiate the effectiveness of the guidance schemes, we present plug-and-play diffusion guidance, namely Self-Attention Guidance (SAG), that can drastically boost the performance of existing diffusion models. Our method, SAG, extracts the intermediate attention map from a diffusion model at every iteration and selects tokens above a certain attention score for masking and blurring to obtain a partially blurred input. Subsequently, we measure the dissimilarity between the predicted noises obtained from feeding the blurred and original input to the diffusion model and leverage it as guidance. With this guidance, we observe apparent improvements in a wide range of diffusion models, e.g., ADM, IDDPM, and Stable Diffusion, and show that the results further improve by combining our method with the conventional guidance scheme. We provide extensive ablation studies to verify our choices.*
|
||||
|
||||
Resources:
|
||||
|
||||
* [Project Page](https://ku-cvlab.github.io/Self-Attention-Guidance).
|
||||
* [Paper](https://arxiv.org/abs/2210.00939).
|
||||
* [Original Code](https://github.com/KU-CVLAB/Self-Attention-Guidance).
|
||||
* [Demo](https://colab.research.google.com/github/SusungHong/Self-Attention-Guidance/blob/main/SAG_Stable.ipynb).
|
||||
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Demo
|
||||
|---|---|:---:|
|
||||
| [StableDiffusionSAGPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_sag.py) | *Text-to-Image Generation* | [Colab](https://colab.research.google.com/github/SusungHong/Self-Attention-Guidance/blob/main/SAG_Stable.ipynb) |
|
||||
|
||||
## Usage example
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import StableDiffusionSAGPipeline
|
||||
from accelerate.utils import set_seed
|
||||
|
||||
pipe = StableDiffusionSAGPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
seed = 8978
|
||||
prompt = "."
|
||||
guidance_scale = 7.5
|
||||
num_images_per_prompt = 1
|
||||
|
||||
sag_scale = 1.0
|
||||
|
||||
set_seed(seed)
|
||||
images = pipe(
|
||||
prompt, num_images_per_prompt=num_images_per_prompt, guidance_scale=guidance_scale, sag_scale=sag_scale
|
||||
).images
|
||||
images[0].save("example.png")
|
||||
```
|
||||
|
||||
## StableDiffusionSAGPipeline
|
||||
[[autodoc]] StableDiffusionSAGPipeline
|
||||
- __call__
|
||||
- all
|
||||
@@ -1,39 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Text-to-Image Generation
|
||||
|
||||
## StableDiffusionPipeline
|
||||
|
||||
The Stable Diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionPipeline`] is capable of generating photo-realistic images given any text input using Stable Diffusion.
|
||||
|
||||
The original codebase can be found here:
|
||||
- *Stable Diffusion V1*: [CompVis/stable-diffusion](https://github.com/CompVis/stable-diffusion)
|
||||
- *Stable Diffusion v2*: [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion)
|
||||
|
||||
Available Checkpoints are:
|
||||
- *stable-diffusion-v1-4 (512x512 resolution)* [CompVis/stable-diffusion-v1-4](https://huggingface.co/CompVis/stable-diffusion-v1-4)
|
||||
- *stable-diffusion-v1-5 (512x512 resolution)* [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5)
|
||||
- *stable-diffusion-2-base (512x512 resolution)*: [stabilityai/stable-diffusion-2-base](https://huggingface.co/stabilityai/stable-diffusion-2-base)
|
||||
- *stable-diffusion-2 (768x768 resolution)*: [stabilityai/stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2)
|
||||
- *stable-diffusion-2-1-base (512x512 resolution)* [stabilityai/stable-diffusion-2-1-base](https://huggingface.co/stabilityai/stable-diffusion-2-1-base)
|
||||
- *stable-diffusion-2-1 (768x768 resolution)*: [stabilityai/stable-diffusion-2-1](https://huggingface.co/stabilityai/stable-diffusion-2-1)
|
||||
|
||||
[[autodoc]] StableDiffusionPipeline
|
||||
- all
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_vae_slicing
|
||||
- disable_vae_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
@@ -1,32 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Super-Resolution
|
||||
|
||||
## StableDiffusionUpscalePipeline
|
||||
|
||||
The upscaler diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), and [LAION](https://laion.ai/), as part of Stable Diffusion 2.0. [`StableDiffusionUpscalePipeline`] can be used to enhance the resolution of input images by a factor of 4.
|
||||
|
||||
The original codebase can be found here:
|
||||
- *Stable Diffusion v2*: [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion#image-upscaling-with-stable-diffusion)
|
||||
|
||||
Available Checkpoints are:
|
||||
- *stabilityai/stable-diffusion-x4-upscaler (x4 resolution resolution)*: [stable-diffusion-x4-upscaler](https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler)
|
||||
|
||||
|
||||
[[autodoc]] StableDiffusionUpscalePipeline
|
||||
- all
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
@@ -1,176 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Stable diffusion 2
|
||||
|
||||
Stable Diffusion 2 is a text-to-image _latent diffusion_ model built upon the work of [Stable Diffusion 1](https://stability.ai/blog/stable-diffusion-public-release).
|
||||
The project to train Stable Diffusion 2 was led by Robin Rombach and Katherine Crowson from [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/).
|
||||
|
||||
*The Stable Diffusion 2.0 release includes robust text-to-image models trained using a brand new text encoder (OpenCLIP), developed by LAION with support from Stability AI, which greatly improves the quality of the generated images compared to earlier V1 releases. The text-to-image models in this release can generate images with default resolutions of both 512x512 pixels and 768x768 pixels.
|
||||
These models are trained on an aesthetic subset of the [LAION-5B dataset](https://laion.ai/blog/laion-5b/) created by the DeepFloyd team at Stability AI, which is then further filtered to remove adult content using [LAION’s NSFW filter](https://openreview.net/forum?id=M3Y74vmsMcY).*
|
||||
|
||||
For more details about how Stable Diffusion 2 works and how it differs from Stable Diffusion 1, please refer to the official [launch announcement post](https://stability.ai/blog/stable-diffusion-v2-release).
|
||||
|
||||
## Tips
|
||||
|
||||
### Available checkpoints:
|
||||
|
||||
Note that the architecture is more or less identical to [Stable Diffusion 1](./stable_diffusion/overview) so please refer to [this page](./stable_diffusion/overview) for API documentation.
|
||||
|
||||
- *Text-to-Image (512x512 resolution)*: [stabilityai/stable-diffusion-2-base](https://huggingface.co/stabilityai/stable-diffusion-2-base) with [`StableDiffusionPipeline`]
|
||||
- *Text-to-Image (768x768 resolution)*: [stabilityai/stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2) with [`StableDiffusionPipeline`]
|
||||
- *Image Inpainting (512x512 resolution)*: [stabilityai/stable-diffusion-2-inpainting](https://huggingface.co/stabilityai/stable-diffusion-2-inpainting) with [`StableDiffusionInpaintPipeline`]
|
||||
- *Super-Resolution (x4 resolution resolution)*: [stable-diffusion-x4-upscaler](https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler) [`StableDiffusionUpscalePipeline`]
|
||||
- *Depth-to-Image (512x512 resolution)*: [stabilityai/stable-diffusion-2-depth](https://huggingface.co/stabilityai/stable-diffusion-2-depth) with [`StableDiffusionDepth2ImagePipeline`]
|
||||
|
||||
We recommend using the [`DPMSolverMultistepScheduler`] as it's currently the fastest scheduler there is.
|
||||
|
||||
|
||||
### Text-to-Image
|
||||
|
||||
- *Text-to-Image (512x512 resolution)*: [stabilityai/stable-diffusion-2-base](https://huggingface.co/stabilityai/stable-diffusion-2-base) with [`StableDiffusionPipeline`]
|
||||
|
||||
```python
|
||||
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
|
||||
import torch
|
||||
|
||||
repo_id = "stabilityai/stable-diffusion-2-base"
|
||||
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
|
||||
|
||||
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "High quality photo of an astronaut riding a horse in space"
|
||||
image = pipe(prompt, num_inference_steps=25).images[0]
|
||||
image.save("astronaut.png")
|
||||
```
|
||||
|
||||
- *Text-to-Image (768x768 resolution)*: [stabilityai/stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2) with [`StableDiffusionPipeline`]
|
||||
|
||||
```python
|
||||
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
|
||||
import torch
|
||||
|
||||
repo_id = "stabilityai/stable-diffusion-2"
|
||||
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
|
||||
|
||||
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "High quality photo of an astronaut riding a horse in space"
|
||||
image = pipe(prompt, guidance_scale=9, num_inference_steps=25).images[0]
|
||||
image.save("astronaut.png")
|
||||
```
|
||||
|
||||
### Image Inpainting
|
||||
|
||||
- *Image Inpainting (512x512 resolution)*: [stabilityai/stable-diffusion-2-inpainting](https://huggingface.co/stabilityai/stable-diffusion-2-inpainting) with [`StableDiffusionInpaintPipeline`]
|
||||
|
||||
```python
|
||||
import PIL
|
||||
import requests
|
||||
import torch
|
||||
from io import BytesIO
|
||||
|
||||
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
|
||||
|
||||
|
||||
def download_image(url):
|
||||
response = requests.get(url)
|
||||
return PIL.Image.open(BytesIO(response.content)).convert("RGB")
|
||||
|
||||
|
||||
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
|
||||
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
|
||||
|
||||
init_image = download_image(img_url).resize((512, 512))
|
||||
mask_image = download_image(mask_url).resize((512, 512))
|
||||
|
||||
repo_id = "stabilityai/stable-diffusion-2-inpainting"
|
||||
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
|
||||
|
||||
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "Face of a yellow cat, high resolution, sitting on a park bench"
|
||||
image = pipe(prompt=prompt, image=init_image, mask_image=mask_image, num_inference_steps=25).images[0]
|
||||
|
||||
image.save("yellow_cat.png")
|
||||
```
|
||||
|
||||
### Super-Resolution
|
||||
|
||||
- *Image Upscaling (x4 resolution resolution)*: [stable-diffusion-x4-upscaler](https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler) with [`StableDiffusionUpscalePipeline`]
|
||||
|
||||
|
||||
```python
|
||||
import requests
|
||||
from PIL import Image
|
||||
from io import BytesIO
|
||||
from diffusers import StableDiffusionUpscalePipeline
|
||||
import torch
|
||||
|
||||
# load model and scheduler
|
||||
model_id = "stabilityai/stable-diffusion-x4-upscaler"
|
||||
pipeline = StableDiffusionUpscalePipeline.from_pretrained(model_id, torch_dtype=torch.float16)
|
||||
pipeline = pipeline.to("cuda")
|
||||
|
||||
# let's download an image
|
||||
url = "https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd2-upscale/low_res_cat.png"
|
||||
response = requests.get(url)
|
||||
low_res_img = Image.open(BytesIO(response.content)).convert("RGB")
|
||||
low_res_img = low_res_img.resize((128, 128))
|
||||
prompt = "a white cat"
|
||||
upscaled_image = pipeline(prompt=prompt, image=low_res_img).images[0]
|
||||
upscaled_image.save("upsampled_cat.png")
|
||||
```
|
||||
|
||||
### Depth-to-Image
|
||||
|
||||
- *Depth-Guided Text-to-Image*: [stabilityai/stable-diffusion-2-depth](https://huggingface.co/stabilityai/stable-diffusion-2-depth) [`StableDiffusionDepth2ImagePipeline`]
|
||||
|
||||
|
||||
```python
|
||||
import torch
|
||||
import requests
|
||||
from PIL import Image
|
||||
|
||||
from diffusers import StableDiffusionDepth2ImgPipeline
|
||||
|
||||
pipe = StableDiffusionDepth2ImgPipeline.from_pretrained(
|
||||
"stabilityai/stable-diffusion-2-depth",
|
||||
torch_dtype=torch.float16,
|
||||
).to("cuda")
|
||||
|
||||
|
||||
url = "http://images.cocodataset.org/val2017/000000039769.jpg"
|
||||
init_image = Image.open(requests.get(url, stream=True).raw)
|
||||
prompt = "two tigers"
|
||||
n_propmt = "bad, deformed, ugly, bad anotomy"
|
||||
image = pipe(prompt=prompt, image=init_image, negative_prompt=n_propmt, strength=0.7).images[0]
|
||||
```
|
||||
|
||||
### How to load and use different schedulers.
|
||||
|
||||
The stable diffusion pipeline uses [`DDIMScheduler`] scheduler by default. But `diffusers` provides many other schedulers that can be used with the stable diffusion pipeline such as [`PNDMScheduler`], [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`] etc.
|
||||
To use a different scheduler, you can either change it via the [`ConfigMixin.from_config`] method or pass the `scheduler` argument to the `from_pretrained` method of the pipeline. For example, to use the [`EulerDiscreteScheduler`], you can do the following:
|
||||
|
||||
```python
|
||||
>>> from diffusers import StableDiffusionPipeline, EulerDiscreteScheduler
|
||||
|
||||
>>> pipeline = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2")
|
||||
>>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
|
||||
|
||||
>>> # or
|
||||
>>> euler_scheduler = EulerDiscreteScheduler.from_pretrained("stabilityai/stable-diffusion-2", subfolder="scheduler")
|
||||
>>> pipeline = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2", scheduler=euler_scheduler)
|
||||
```
|
||||
@@ -1,90 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Safe Stable Diffusion
|
||||
|
||||
Safe Stable Diffusion was proposed in [Safe Latent Diffusion: Mitigating Inappropriate Degeneration in Diffusion Models](https://arxiv.org/abs/2211.05105) and mitigates the well known issue that models like Stable Diffusion that are trained on unfiltered, web-crawled datasets tend to suffer from inappropriate degeneration. For instance Stable Diffusion may unexpectedly generate nudity, violence, images depicting self-harm, or otherwise offensive content.
|
||||
Safe Stable Diffusion is an extension to the Stable Diffusion that drastically reduces content like this.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*Text-conditioned image generation models have recently achieved astonishing results in image quality and text alignment and are consequently employed in a fast-growing number of applications. Since they are highly data-driven, relying on billion-sized datasets randomly scraped from the internet, they also suffer, as we demonstrate, from degenerated and biased human behavior. In turn, they may even reinforce such biases. To help combat these undesired side effects, we present safe latent diffusion (SLD). Specifically, to measure the inappropriate degeneration due to unfiltered and imbalanced training sets, we establish a novel image generation test bed-inappropriate image prompts (I2P)-containing dedicated, real-world image-to-text prompts covering concepts such as nudity and violence. As our exhaustive empirical evaluation demonstrates, the introduced SLD removes and suppresses inappropriate image parts during the diffusion process, with no additional training required and no adverse effect on overall image quality or text alignment.*
|
||||
|
||||
|
||||
*Overview*:
|
||||
|
||||
| Pipeline | Tasks | Colab | Demo
|
||||
|---|---|:---:|:---:|
|
||||
| [pipeline_stable_diffusion_safe.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion_safe/pipeline_stable_diffusion_safe.py) | *Text-to-Image Generation* | [](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb) | [](https://huggingface.co/spaces/AIML-TUDA/unsafe-vs-safe-stable-diffusion)
|
||||
|
||||
## Tips
|
||||
|
||||
- Safe Stable Diffusion may also be used with weights of [Stable Diffusion](./api/pipelines/stable_diffusion/text2img).
|
||||
|
||||
### Run Safe Stable Diffusion
|
||||
|
||||
Safe Stable Diffusion can be tested very easily with the [`StableDiffusionPipelineSafe`], and the `"AIML-TUDA/stable-diffusion-safe"` checkpoint exactly in the same way it is shown in the [Conditional Image Generation Guide](./using-diffusers/conditional_image_generation).
|
||||
|
||||
### Interacting with the Safety Concept
|
||||
|
||||
To check and edit the currently used safety concept, use the `safety_concept` property of [`StableDiffusionPipelineSafe`]
|
||||
```python
|
||||
>>> from diffusers import StableDiffusionPipelineSafe
|
||||
|
||||
>>> pipeline = StableDiffusionPipelineSafe.from_pretrained("AIML-TUDA/stable-diffusion-safe")
|
||||
>>> pipeline.safety_concept
|
||||
```
|
||||
For each image generation the active concept is also contained in [`StableDiffusionSafePipelineOutput`].
|
||||
|
||||
### Using pre-defined safety configurations
|
||||
|
||||
You may use the 4 configurations defined in the [Safe Latent Diffusion paper](https://arxiv.org/abs/2211.05105) as follows:
|
||||
|
||||
```python
|
||||
>>> from diffusers import StableDiffusionPipelineSafe
|
||||
>>> from diffusers.pipelines.stable_diffusion_safe import SafetyConfig
|
||||
|
||||
>>> pipeline = StableDiffusionPipelineSafe.from_pretrained("AIML-TUDA/stable-diffusion-safe")
|
||||
>>> prompt = "the four horsewomen of the apocalypse, painting by tom of finland, gaston bussiere, craig mullins, j. c. leyendecker"
|
||||
>>> out = pipeline(prompt=prompt, **SafetyConfig.MAX)
|
||||
```
|
||||
|
||||
The following configurations are available: `SafetyConfig.WEAK`, `SafetyConfig.MEDIUM`, `SafetyConfig.STRONG`, and `SafetyConfig.MAX`.
|
||||
|
||||
### How to load and use different schedulers.
|
||||
|
||||
The safe stable diffusion pipeline uses [`PNDMScheduler`] scheduler by default. But `diffusers` provides many other schedulers that can be used with the stable diffusion pipeline such as [`DDIMScheduler`], [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`] etc.
|
||||
To use a different scheduler, you can either change it via the [`ConfigMixin.from_config`] method or pass the `scheduler` argument to the `from_pretrained` method of the pipeline. For example, to use the [`EulerDiscreteScheduler`], you can do the following:
|
||||
|
||||
```python
|
||||
>>> from diffusers import StableDiffusionPipelineSafe, EulerDiscreteScheduler
|
||||
|
||||
>>> pipeline = StableDiffusionPipelineSafe.from_pretrained("AIML-TUDA/stable-diffusion-safe")
|
||||
>>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
|
||||
|
||||
>>> # or
|
||||
>>> euler_scheduler = EulerDiscreteScheduler.from_pretrained("AIML-TUDA/stable-diffusion-safe", subfolder="scheduler")
|
||||
>>> pipeline = StableDiffusionPipelineSafe.from_pretrained(
|
||||
... "AIML-TUDA/stable-diffusion-safe", scheduler=euler_scheduler
|
||||
... )
|
||||
```
|
||||
|
||||
|
||||
## StableDiffusionSafePipelineOutput
|
||||
[[autodoc]] pipelines.stable_diffusion_safe.StableDiffusionSafePipelineOutput
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## StableDiffusionPipelineSafe
|
||||
[[autodoc]] StableDiffusionPipelineSafe
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,97 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Stable unCLIP
|
||||
|
||||
Stable unCLIP checkpoints are finetuned from [stable diffusion 2.1](./stable_diffusion_2) checkpoints to condition on CLIP image embeddings.
|
||||
Stable unCLIP also still conditions on text embeddings. Given the two separate conditionings, stable unCLIP can be used
|
||||
for text guided image variation. When combined with an unCLIP prior, it can also be used for full text to image generation.
|
||||
|
||||
## Tips
|
||||
|
||||
Stable unCLIP takes a `noise_level` as input during inference. `noise_level` determines how much noise is added
|
||||
to the image embeddings. A higher `noise_level` increases variation in the final un-noised images. By default,
|
||||
we do not add any additional noise to the image embeddings i.e. `noise_level = 0`.
|
||||
|
||||
### Available checkpoints:
|
||||
|
||||
TODO
|
||||
|
||||
### Text-to-Image Generation
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import StableUnCLIPPipeline
|
||||
|
||||
pipe = StableUnCLIPPipeline.from_pretrained(
|
||||
"fusing/stable-unclip-2-1-l", torch_dtype=torch.float16
|
||||
) # TODO update model path
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "a photo of an astronaut riding a horse on mars"
|
||||
images = pipe(prompt).images
|
||||
images[0].save("astronaut_horse.png")
|
||||
```
|
||||
|
||||
|
||||
### Text guided Image-to-Image Variation
|
||||
|
||||
```python
|
||||
import requests
|
||||
import torch
|
||||
from PIL import Image
|
||||
from io import BytesIO
|
||||
|
||||
from diffusers import StableUnCLIPImg2ImgPipeline
|
||||
|
||||
pipe = StableUnCLIPImg2ImgPipeline.from_pretrained(
|
||||
"fusing/stable-unclip-2-1-l-img2img", torch_dtype=torch.float16
|
||||
) # TODO update model path
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
|
||||
|
||||
response = requests.get(url)
|
||||
init_image = Image.open(BytesIO(response.content)).convert("RGB")
|
||||
init_image = init_image.resize((768, 512))
|
||||
|
||||
prompt = "A fantasy landscape, trending on artstation"
|
||||
|
||||
images = pipe(prompt, init_image).images
|
||||
images[0].save("fantasy_landscape.png")
|
||||
```
|
||||
|
||||
### StableUnCLIPPipeline
|
||||
|
||||
[[autodoc]] StableUnCLIPPipeline
|
||||
- all
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_vae_slicing
|
||||
- disable_vae_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
|
||||
|
||||
### StableUnCLIPImg2ImgPipeline
|
||||
|
||||
[[autodoc]] StableUnCLIPImg2ImgPipeline
|
||||
- all
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_vae_slicing
|
||||
- disable_vae_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
|
||||
@@ -1,37 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# unCLIP
|
||||
|
||||
## Overview
|
||||
|
||||
[Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) by Aditya Ramesh, Prafulla Dhariwal, Alex Nichol, Casey Chu, Mark Chen
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
Contrastive models like CLIP have been shown to learn robust representations of images that capture both semantics and style. To leverage these representations for image generation, we propose a two-stage model: a prior that generates a CLIP image embedding given a text caption, and a decoder that generates an image conditioned on the image embedding. We show that explicitly generating image representations improves image diversity with minimal loss in photorealism and caption similarity. Our decoders conditioned on image representations can also produce variations of an image that preserve both its semantics and style, while varying the non-essential details absent from the image representation. Moreover, the joint embedding space of CLIP enables language-guided image manipulations in a zero-shot fashion. We use diffusion models for the decoder and experiment with both autoregressive and diffusion models for the prior, finding that the latter are computationally more efficient and produce higher-quality samples.
|
||||
|
||||
The unCLIP model in diffusers comes from kakaobrain's karlo and the original codebase can be found [here](https://github.com/kakaobrain/karlo). Additionally, lucidrains has a DALL-E 2 recreation [here](https://github.com/lucidrains/DALLE2-pytorch).
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Colab
|
||||
|---|---|:---:|
|
||||
| [pipeline_unclip.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/unclip/pipeline_unclip.py) | *Text-to-Image Generation* | - |
|
||||
| [pipeline_unclip_image_variation.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/unclip/pipeline_unclip_image_variation.py) | *Image-Guided Image Generation* | - |
|
||||
|
||||
|
||||
## UnCLIPPipeline
|
||||
[[autodoc]] UnCLIPPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
[[autodoc]] UnCLIPImageVariationPipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,70 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# VersatileDiffusion
|
||||
|
||||
VersatileDiffusion was proposed in [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) by Xingqian Xu, Zhangyang Wang, Eric Zhang, Kai Wang, Humphrey Shi .
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
*The recent advances in diffusion models have set an impressive milestone in many generation tasks. Trending works such as DALL-E2, Imagen, and Stable Diffusion have attracted great interest in academia and industry. Despite the rapid landscape changes, recent new approaches focus on extensions and performance rather than capacity, thus requiring separate models for separate tasks. In this work, we expand the existing single-flow diffusion pipeline into a multi-flow network, dubbed Versatile Diffusion (VD), that handles text-to-image, image-to-text, image-variation, and text-variation in one unified model. Moreover, we generalize VD to a unified multi-flow multimodal diffusion framework with grouped layers, swappable streams, and other propositions that can process modalities beyond images and text. Through our experiments, we demonstrate that VD and its underlying framework have the following merits: a) VD handles all subtasks with competitive quality; b) VD initiates novel extensions and applications such as disentanglement of style and semantic, image-text dual-guided generation, etc.; c) Through these experiments and applications, VD provides more semantic insights of the generated outputs.*
|
||||
|
||||
## Tips
|
||||
|
||||
- VersatileDiffusion is conceptually very similar as [Stable Diffusion](./api/pipelines/stable_diffusion/overview), but instead of providing just a image data stream conditioned on text, VersatileDiffusion provides both a image and text data stream and can be conditioned on both text and image.
|
||||
|
||||
### *Run VersatileDiffusion*
|
||||
|
||||
You can both load the memory intensive "all-in-one" [`VersatileDiffusionPipeline`] that can run all tasks
|
||||
with the same class as shown in [`VersatileDiffusionPipeline.text_to_image`], [`VersatileDiffusionPipeline.image_variation`], and [`VersatileDiffusionPipeline.dual_guided`]
|
||||
|
||||
**or**
|
||||
|
||||
You can run the individual pipelines which are much more memory efficient:
|
||||
|
||||
- *Text-to-Image*: [`VersatileDiffusionTextToImagePipeline.__call__`]
|
||||
- *Image Variation*: [`VersatileDiffusionImageVariationPipeline.__call__`]
|
||||
- *Dual Text and Image Guided Generation*: [`VersatileDiffusionDualGuidedPipeline.__call__`]
|
||||
|
||||
### *How to load and use different schedulers.*
|
||||
|
||||
The versatile diffusion pipelines uses [`DDIMScheduler`] scheduler by default. But `diffusers` provides many other schedulers that can be used with the alt diffusion pipeline such as [`PNDMScheduler`], [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`] etc.
|
||||
To use a different scheduler, you can either change it via the [`ConfigMixin.from_config`] method or pass the `scheduler` argument to the `from_pretrained` method of the pipeline. For example, to use the [`EulerDiscreteScheduler`], you can do the following:
|
||||
|
||||
```python
|
||||
>>> from diffusers import VersatileDiffusionPipeline, EulerDiscreteScheduler
|
||||
|
||||
>>> pipeline = VersatileDiffusionPipeline.from_pretrained("shi-labs/versatile-diffusion")
|
||||
>>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
|
||||
|
||||
>>> # or
|
||||
>>> euler_scheduler = EulerDiscreteScheduler.from_pretrained("shi-labs/versatile-diffusion", subfolder="scheduler")
|
||||
>>> pipeline = VersatileDiffusionPipeline.from_pretrained("shi-labs/versatile-diffusion", scheduler=euler_scheduler)
|
||||
```
|
||||
|
||||
## VersatileDiffusionPipeline
|
||||
[[autodoc]] VersatileDiffusionPipeline
|
||||
|
||||
## VersatileDiffusionTextToImagePipeline
|
||||
[[autodoc]] VersatileDiffusionTextToImagePipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## VersatileDiffusionImageVariationPipeline
|
||||
[[autodoc]] VersatileDiffusionImageVariationPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## VersatileDiffusionDualGuidedPipeline
|
||||
[[autodoc]] VersatileDiffusionDualGuidedPipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,35 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# VQDiffusion
|
||||
|
||||
## Overview
|
||||
|
||||
[Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) by Shuyang Gu, Dong Chen, Jianmin Bao, Fang Wen, Bo Zhang, Dongdong Chen, Lu Yuan, Baining Guo
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
We present the vector quantized diffusion (VQ-Diffusion) model for text-to-image generation. This method is based on a vector quantized variational autoencoder (VQ-VAE) whose latent space is modeled by a conditional variant of the recently developed Denoising Diffusion Probabilistic Model (DDPM). We find that this latent-space method is well-suited for text-to-image generation tasks because it not only eliminates the unidirectional bias with existing methods but also allows us to incorporate a mask-and-replace diffusion strategy to avoid the accumulation of errors, which is a serious problem with existing methods. Our experiments show that the VQ-Diffusion produces significantly better text-to-image generation results when compared with conventional autoregressive (AR) models with similar numbers of parameters. Compared with previous GAN-based text-to-image methods, our VQ-Diffusion can handle more complex scenes and improve the synthesized image quality by a large margin. Finally, we show that the image generation computation in our method can be made highly efficient by reparameterization. With traditional AR methods, the text-to-image generation time increases linearly with the output image resolution and hence is quite time consuming even for normal size images. The VQ-Diffusion allows us to achieve a better trade-off between quality and speed. Our experiments indicate that the VQ-Diffusion model with the reparameterization is fifteen times faster than traditional AR methods while achieving a better image quality.
|
||||
|
||||
The original codebase can be found [here](https://github.com/microsoft/VQ-Diffusion).
|
||||
|
||||
## Available Pipelines:
|
||||
|
||||
| Pipeline | Tasks | Colab
|
||||
|---|---|:---:|
|
||||
| [pipeline_vq_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/vq_diffusion/pipeline_vq_diffusion.py) | *Text-to-Image Generation* | - |
|
||||
|
||||
|
||||
## VQDiffusionPipeline
|
||||
[[autodoc]] VQDiffusionPipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,27 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Denoising diffusion implicit models (DDIM)
|
||||
|
||||
## Overview
|
||||
|
||||
[Denoising Diffusion Implicit Models](https://arxiv.org/abs/2010.02502) (DDIM) by Jiaming Song, Chenlin Meng and Stefano Ermon.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
Denoising diffusion probabilistic models (DDPMs) have achieved high quality image generation without adversarial training, yet they require simulating a Markov chain for many steps to produce a sample. To accelerate sampling, we present denoising diffusion implicit models (DDIMs), a more efficient class of iterative implicit probabilistic models with the same training procedure as DDPMs. In DDPMs, the generative process is defined as the reverse of a Markovian diffusion process. We construct a class of non-Markovian diffusion processes that lead to the same training objective, but whose reverse process can be much faster to sample from. We empirically demonstrate that DDIMs can produce high quality samples 10× to 50× faster in terms of wall-clock time compared to DDPMs, allow us to trade off computation for sample quality, and can perform semantically meaningful image interpolation directly in the latent space.
|
||||
|
||||
The original codebase of this paper can be found here: [ermongroup/ddim](https://github.com/ermongroup/ddim).
|
||||
For questions, feel free to contact the author on [tsong.me](https://tsong.me/).
|
||||
|
||||
## DDIMScheduler
|
||||
[[autodoc]] DDIMScheduler
|
||||
@@ -1,21 +0,0 @@
|
||||
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Inverse Denoising Diffusion Implicit Models (DDIMInverse)
|
||||
|
||||
## Overview
|
||||
|
||||
This scheduler is the inverted scheduler of [Denoising Diffusion Implicit Models](https://arxiv.org/abs/2010.02502) (DDIM) by Jiaming Song, Chenlin Meng and Stefano Ermon.
|
||||
The implementation is mostly based on the DDIM inversion definition of [Null-text Inversion for Editing Real Images using Guided Diffusion Models](https://arxiv.org/pdf/2211.09794.pdf)
|
||||
|
||||
## DDIMInverseScheduler
|
||||
[[autodoc]] DDIMInverseScheduler
|
||||
@@ -1,27 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Denoising diffusion probabilistic models (DDPM)
|
||||
|
||||
## Overview
|
||||
|
||||
[Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2006.11239)
|
||||
(DDPM) by Jonathan Ho, Ajay Jain and Pieter Abbeel proposes the diffusion based model of the same name, but in the context of the 🤗 Diffusers library, DDPM refers to the discrete denoising scheduler from the paper as well as the pipeline.
|
||||
|
||||
The abstract of the paper is the following:
|
||||
|
||||
We present high quality image synthesis results using diffusion probabilistic models, a class of latent variable models inspired by considerations from nonequilibrium thermodynamics. Our best results are obtained by training on a weighted variational bound designed according to a novel connection between diffusion probabilistic models and denoising score matching with Langevin dynamics, and our models naturally admit a progressive lossy decompression scheme that can be interpreted as a generalization of autoregressive decoding. On the unconditional CIFAR10 dataset, we obtain an Inception score of 9.46 and a state-of-the-art FID score of 3.17. On 256x256 LSUN, we obtain sample quality similar to ProgressiveGAN.
|
||||
|
||||
The original paper can be found [here](https://arxiv.org/abs/2010.02502).
|
||||
|
||||
## DDPMScheduler
|
||||
[[autodoc]] DDPMScheduler
|
||||
@@ -1,22 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# DEIS
|
||||
|
||||
Fast Sampling of Diffusion Models with Exponential Integrator.
|
||||
|
||||
## Overview
|
||||
|
||||
Original paper can be found [here](https://arxiv.org/abs/2204.13902). The original implementation can be found [here](https://github.com/qsh-zh/deis).
|
||||
|
||||
## DEISMultistepScheduler
|
||||
[[autodoc]] DEISMultistepScheduler
|
||||
@@ -1,22 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# DPM Discrete Scheduler inspired by Karras et. al paper
|
||||
|
||||
## Overview
|
||||
|
||||
Inspired by [Karras et. al](https://arxiv.org/abs/2206.00364). Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
|
||||
|
||||
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
|
||||
|
||||
## KDPM2DiscreteScheduler
|
||||
[[autodoc]] KDPM2DiscreteScheduler
|
||||
@@ -1,22 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# DPM Discrete Scheduler with ancestral sampling inspired by Karras et. al paper
|
||||
|
||||
## Overview
|
||||
|
||||
Inspired by [Karras et. al](https://arxiv.org/abs/2206.00364). Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
|
||||
|
||||
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
|
||||
|
||||
## KDPM2AncestralDiscreteScheduler
|
||||
[[autodoc]] KDPM2AncestralDiscreteScheduler
|
||||
@@ -1,21 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Euler scheduler
|
||||
|
||||
## Overview
|
||||
|
||||
Euler scheduler (Algorithm 2) from the paper [Elucidating the Design Space of Diffusion-Based Generative Models](https://arxiv.org/abs/2206.00364) by Karras et al. (2022). Based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L51) implementation by Katherine Crowson.
|
||||
Fast scheduler which often times generates good outputs with 20-30 steps.
|
||||
|
||||
## EulerDiscreteScheduler
|
||||
[[autodoc]] EulerDiscreteScheduler
|
||||
@@ -1,21 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Euler Ancestral scheduler
|
||||
|
||||
## Overview
|
||||
|
||||
Ancestral sampling with Euler method steps. Based on the original (k-diffusion)[https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72] implementation by Katherine Crowson.
|
||||
Fast scheduler which often times generates good outputs with 20-30 steps.
|
||||
|
||||
## EulerAncestralDiscreteScheduler
|
||||
[[autodoc]] EulerAncestralDiscreteScheduler
|
||||
@@ -1,23 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Heun scheduler inspired by Karras et. al paper
|
||||
|
||||
## Overview
|
||||
|
||||
Algorithm 1 of [Karras et. al](https://arxiv.org/abs/2206.00364).
|
||||
Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
|
||||
|
||||
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
|
||||
|
||||
## HeunDiscreteScheduler
|
||||
[[autodoc]] HeunDiscreteScheduler
|
||||
@@ -1,20 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# improved pseudo numerical methods for diffusion models (iPNDM)
|
||||
|
||||
## Overview
|
||||
|
||||
Original implementation can be found [here](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296).
|
||||
|
||||
## IPNDMScheduler
|
||||
[[autodoc]] IPNDMScheduler
|
||||
@@ -1,20 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Linear multistep scheduler for discrete beta schedules
|
||||
|
||||
## Overview
|
||||
|
||||
Original implementation can be found [here](https://arxiv.org/abs/2206.00364).
|
||||
|
||||
## LMSDiscreteScheduler
|
||||
[[autodoc]] LMSDiscreteScheduler
|
||||
@@ -1,20 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Multistep DPM-Solver
|
||||
|
||||
## Overview
|
||||
|
||||
Original paper can be found [here](https://arxiv.org/abs/2206.00927) and the [improved version](https://arxiv.org/abs/2211.01095). The original implementation can be found [here](https://github.com/LuChengTHU/dpm-solver).
|
||||
|
||||
## DPMSolverMultistepScheduler
|
||||
[[autodoc]] DPMSolverMultistepScheduler
|
||||
@@ -1,92 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Schedulers
|
||||
|
||||
Diffusers contains multiple pre-built schedule functions for the diffusion process.
|
||||
|
||||
## What is a scheduler?
|
||||
|
||||
The schedule functions, denoted *Schedulers* in the library take in the output of a trained model, a sample which the diffusion process is iterating on, and a timestep to return a denoised sample. That's why schedulers may also be called *Samplers* in other diffusion models implementations.
|
||||
|
||||
- Schedulers define the methodology for iteratively adding noise to an image or for updating a sample based on model outputs.
|
||||
- adding noise in different manners represent the algorithmic processes to train a diffusion model by adding noise to images.
|
||||
- for inference, the scheduler defines how to update a sample based on an output from a pretrained model.
|
||||
- Schedulers are often defined by a *noise schedule* and an *update rule* to solve the differential equation solution.
|
||||
|
||||
### Discrete versus continuous schedulers
|
||||
|
||||
All schedulers take in a timestep to predict the updated version of the sample being diffused.
|
||||
The timesteps dictate where in the diffusion process the step is, where data is generated by iterating forward in time and inference is executed by propagating backwards through timesteps.
|
||||
Different algorithms use timesteps that can be discrete (accepting `int` inputs), such as the [`DDPMScheduler`] or [`PNDMScheduler`], or continuous (accepting `float` inputs), such as the score-based schedulers [`ScoreSdeVeScheduler`] or [`ScoreSdeVpScheduler`].
|
||||
|
||||
## Designing Re-usable schedulers
|
||||
|
||||
The core design principle between the schedule functions is to be model, system, and framework independent.
|
||||
This allows for rapid experimentation and cleaner abstractions in the code, where the model prediction is separated from the sample update.
|
||||
To this end, the design of schedulers is such that:
|
||||
|
||||
- Schedulers can be used interchangeably between diffusion models in inference to find the preferred trade-off between speed and generation quality.
|
||||
- Schedulers are currently by default in PyTorch, but are designed to be framework independent (partial Jax support currently exists).
|
||||
- Many diffusion pipelines, such as [`StableDiffusionPipeline`] and [`DiTPipeline`] can use any of [`KarrasDiffusionSchedulers`]
|
||||
|
||||
## Schedulers Summary
|
||||
|
||||
The following table summarizes all officially supported schedulers, their corresponding paper
|
||||
|
||||
| Scheduler | Paper |
|
||||
|---|---|
|
||||
| [ddim](./ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) |
|
||||
| [ddim_inverse](./ddim_inverse) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) |
|
||||
| [ddpm](./ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) |
|
||||
| [deis](./deis) | [**DEISMultistepScheduler**](https://arxiv.org/abs/2204.13902) |
|
||||
| [singlestep_dpm_solver](./singlestep_dpm_solver) | [**Singlestep DPM-Solver**](https://arxiv.org/abs/2206.00927) |
|
||||
| [multistep_dpm_solver](./multistep_dpm_solver) | [**Multistep DPM-Solver**](https://arxiv.org/abs/2206.00927) |
|
||||
| [heun](./heun) | [**Heun scheduler inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
|
||||
| [dpm_discrete](./dpm_discrete) | [**DPM Discrete Scheduler inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
|
||||
| [dpm_discrete_ancestral](./dpm_discrete_ancestral) | [**DPM Discrete Scheduler with ancestral sampling inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
|
||||
| [stochastic_karras_ve](./stochastic_karras_ve) | [**Variance exploding, stochastic sampling from Karras et. al**](https://arxiv.org/abs/2206.00364) |
|
||||
| [lms_discrete](./lms_discrete) | [**Linear multistep scheduler for discrete beta schedules**](https://arxiv.org/abs/2206.00364) |
|
||||
| [pndm](./pndm) | [**Pseudo numerical methods for diffusion models (PNDM)**](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181) |
|
||||
| [score_sde_ve](./score_sde_ve) | [**variance exploding stochastic differential equation (VE-SDE) scheduler**](https://arxiv.org/abs/2011.13456) |
|
||||
| [ipndm](./ipndm) | [**improved pseudo numerical methods for diffusion models (iPNDM)**](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296) |
|
||||
| [score_sde_vp](./score_sde_vp) | [**Variance preserving stochastic differential equation (VP-SDE) scheduler**](https://arxiv.org/abs/2011.13456) |
|
||||
| [euler](./euler) | [**Euler scheduler**](https://arxiv.org/abs/2206.00364) |
|
||||
| [euler_ancestral](./euler_ancestral) | [**Euler Ancestral scheduler**](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72) |
|
||||
| [vq_diffusion](./vq_diffusion) | [**VQDiffusionScheduler**](https://arxiv.org/abs/2111.14822) |
|
||||
| [unipc](./unipc) | [**UniPCMultistepScheduler**](https://arxiv.org/abs/2302.04867) |
|
||||
| [repaint](./repaint) | [**RePaint scheduler**](https://arxiv.org/abs/2201.09865) |
|
||||
|
||||
## API
|
||||
|
||||
The core API for any new scheduler must follow a limited structure.
|
||||
- Schedulers should provide one or more `def step(...)` functions that should be called to update the generated sample iteratively.
|
||||
- Schedulers should provide a `set_timesteps(...)` method that configures the parameters of a schedule function for a specific inference task.
|
||||
- Schedulers should be framework-specific.
|
||||
|
||||
The base class [`SchedulerMixin`] implements low level utilities used by multiple schedulers.
|
||||
|
||||
### SchedulerMixin
|
||||
[[autodoc]] SchedulerMixin
|
||||
|
||||
### SchedulerOutput
|
||||
The class [`SchedulerOutput`] contains the outputs from any schedulers `step(...)` call.
|
||||
|
||||
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
|
||||
|
||||
### KarrasDiffusionSchedulers
|
||||
|
||||
`KarrasDiffusionSchedulers` encompasses the main generalization of schedulers in Diffusers. The schedulers in this class are distinguished, at a high level, by their noise sampling strategy; the type of network and scaling; and finally the training strategy or how the loss is weighed.
|
||||
|
||||
The different schedulers, depending on the type of ODE solver, fall into the above taxonomy and provide a good abstraction for the design of the main schedulers implemented in Diffusers. The schedulers in this class are given below:
|
||||
|
||||
[[autodoc]] schedulers.scheduling_utils.KarrasDiffusionSchedulers
|
||||
@@ -1,20 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Pseudo numerical methods for diffusion models (PNDM)
|
||||
|
||||
## Overview
|
||||
|
||||
Original implementation can be found [here](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181).
|
||||
|
||||
## PNDMScheduler
|
||||
[[autodoc]] PNDMScheduler
|
||||
@@ -1,23 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# RePaint scheduler
|
||||
|
||||
## Overview
|
||||
|
||||
DDPM-based inpainting scheduler for unsupervised inpainting with extreme masks.
|
||||
Intended for use with [`RePaintPipeline`].
|
||||
Based on the paper [RePaint: Inpainting using Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2201.09865)
|
||||
and the original implementation by Andreas Lugmayr et al.: https://github.com/andreas128/RePaint
|
||||
|
||||
## RePaintScheduler
|
||||
[[autodoc]] RePaintScheduler
|
||||
@@ -1,20 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# variance exploding stochastic differential equation (VE-SDE) scheduler
|
||||
|
||||
## Overview
|
||||
|
||||
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
|
||||
|
||||
## ScoreSdeVeScheduler
|
||||
[[autodoc]] ScoreSdeVeScheduler
|
||||
@@ -1,20 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Singlestep DPM-Solver
|
||||
|
||||
## Overview
|
||||
|
||||
Original paper can be found [here](https://arxiv.org/abs/2206.00927) and the [improved version](https://arxiv.org/abs/2211.01095). The original implementation can be found [here](https://github.com/LuChengTHU/dpm-solver).
|
||||
|
||||
## DPMSolverSinglestepScheduler
|
||||
[[autodoc]] DPMSolverSinglestepScheduler
|
||||
@@ -1,20 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Variance exploding, stochastic sampling from Karras et. al
|
||||
|
||||
## Overview
|
||||
|
||||
Original paper can be found [here](https://arxiv.org/abs/2206.00364).
|
||||
|
||||
## KarrasVeScheduler
|
||||
[[autodoc]] KarrasVeScheduler
|
||||
@@ -1,24 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# UniPC
|
||||
|
||||
## Overview
|
||||
|
||||
UniPC is a training-free framework designed for the fast sampling of diffusion models, which consists of a corrector (UniC) and a predictor (UniP) that share a unified analytical form and support arbitrary orders.
|
||||
|
||||
For more details about the method, please refer to the [[paper]](https://arxiv.org/abs/2302.04867) and the [[code]](https://github.com/wl-zhao/UniPC).
|
||||
|
||||
Fast Sampling of Diffusion Models with Exponential Integrator.
|
||||
|
||||
## UniPCMultistepScheduler
|
||||
[[autodoc]] UniPCMultistepScheduler
|
||||
@@ -1,20 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# VQDiffusionScheduler
|
||||
|
||||
## Overview
|
||||
|
||||
Original paper can be found [here](https://arxiv.org/abs/2111.14822)
|
||||
|
||||
## VQDiffusionScheduler
|
||||
[[autodoc]] VQDiffusionScheduler
|
||||
@@ -1,49 +0,0 @@
|
||||
# 🧨 Diffusers’ Ethical Guidelines
|
||||
|
||||
## Preamble
|
||||
|
||||
[Diffusers](https://huggingface.co/docs/diffusers/index) provides pre-trained diffusion models and serves as a modular toolbox for inference and training.
|
||||
|
||||
Given its real case applications in the world and potential negative impacts on society, we think it is important to provide the project with ethical guidelines to guide the development, users’ contributions, and usage of the Diffusers library.
|
||||
|
||||
The risks associated with using this technology are still being examined, but to name a few: copyrights issues for artists; deep-fake exploitation; sexual content generation in inappropriate contexts; non-consensual impersonation; harmful social biases perpetuating the oppression of marginalized groups.
|
||||
We will keep tracking risks and adapt the following guidelines based on the community's responsiveness and valuable feedback.
|
||||
|
||||
|
||||
## Scope
|
||||
|
||||
The Diffusers community will apply the following ethical guidelines to the project’s development and help coordinate how the community will integrate the contributions, especially concerning sensitive topics related to ethical concerns.
|
||||
|
||||
|
||||
## Ethical guidelines
|
||||
|
||||
The following ethical guidelines apply generally, but we will primarily implement them when dealing with ethically sensitive issues while making a technical choice. Furthermore, we commit to adapting those ethical principles over time following emerging harms related to the state of the art of the technology in question.
|
||||
|
||||
- **Transparency**: we are committed to being transparent in managing PRs, explaining our choices to users, and making technical decisions.
|
||||
|
||||
- **Consistency**: we are committed to guaranteeing our users the same level of attention in project management, keeping it technically stable and consistent.
|
||||
|
||||
- **Simplicity**: with a desire to make it easy to use and exploit the Diffusers library, we are committed to keeping the project’s goals lean and coherent.
|
||||
|
||||
- **Accessibility**: the Diffusers project helps lower the entry bar for contributors who can help run it even without technical expertise. Doing so makes research artifacts more accessible to the community.
|
||||
|
||||
- **Reproducibility**: we aim to be transparent about the reproducibility of upstream code, models, and datasets when made available through the Diffusers library.
|
||||
|
||||
- **Responsibility**: as a community and through teamwork, we hold a collective responsibility to our users by anticipating and mitigating this technology's potential risks and dangers.
|
||||
|
||||
|
||||
## Examples of implementations: Safety features and Mechanisms
|
||||
|
||||
The team works daily to make the technical and non-technical tools available to deal with the potential ethical and social risks associated with diffusion technology. Moreover, the community's input is invaluable in ensuring these features' implementation and raising awareness with us.
|
||||
|
||||
- [**Community tab**](https://huggingface.co/docs/hub/repositories-pull-requests-discussions): it enables the community to discuss and better collaborate on a project.
|
||||
|
||||
- **Bias exploration and evaluation**: the Hugging Face team provides a [space](https://huggingface.co/spaces/society-ethics/DiffusionBiasExplorer) to demonstrate the biases in Stable Diffusion interactively. In this sense, we support and encourage bias explorers and evaluations.
|
||||
|
||||
- **Encouraging safety in deployment**
|
||||
|
||||
- [**Safe Stable Diffusion**](https://huggingface.co/docs/diffusers/main/en/api/pipelines/stable_diffusion_safe): It mitigates the well-known issue that models, like Stable Diffusion, that are trained on unfiltered, web-crawled datasets tend to suffer from inappropriate degeneration. Related paper: [Safe Latent Diffusion: Mitigating Inappropriate Degeneration in Diffusion Models](https://arxiv.org/abs/2211.05105).
|
||||
|
||||
- **Staged released on the Hub**: in particularly sensitive situations, access to some repositories should be restricted. This staged release is an intermediary step that allows the repository’s authors to have more control over its use.
|
||||
|
||||
- **Licensing**: [OpenRAILs](https://huggingface.co/blog/open_rail), a new type of licensing, allow us to ensure free access while having a set of restrictions that ensure more responsible use.
|
||||
@@ -1,110 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Philosophy
|
||||
|
||||
🧨 Diffusers provides **state-of-the-art** pretrained diffusion models across multiple modalities.
|
||||
Its purpose is to serve as a **modular toolbox** for both inference and training.
|
||||
|
||||
We aim at building a library that stands the test of time and therefore take API design very seriously.
|
||||
|
||||
In a nutshell, Diffusers is built to be a natural extension of PyTorch. Therefore, most of our design choices are based on [PyTorch's Design Principles](https://pytorch.org/docs/stable/community/design.html#pytorch-design-philosophy). Let's go over the most important ones:
|
||||
|
||||
## Usability over Performance
|
||||
|
||||
- While Diffusers has many built-in performance-enhancing features (see [Memory and Speed](https://huggingface.co/docs/diffusers/optimization/fp16)), models are always loaded with the highest precision and lowest optimization. Therefore, by default diffusion pipelines are always instantiated on CPU with float32 precision if not otherwise defined by the user. This ensures usability across different platforms and accelerators and means that no complex installations are required to run the library.
|
||||
- Diffusers aim at being a **light-weight** package and therefore has very few required dependencies, but many soft dependencies that can improve performance (such as `accelerate`, `safetensors`, `onnx`, etc...). We strive to keep the library as lightweight as possible so that it can be added without much concern as a dependency on other packages.
|
||||
- Diffusers prefers simple, self-explainable code over condensed, magic code. This means that short-hand code syntaxes such as lambda functions, and advanced PyTorch operators are often not desired.
|
||||
|
||||
## Simple over easy
|
||||
|
||||
As PyTorch states, **explicit is better than implicit** and **simple is better than complex**. This design philosophy is reflected in multiple parts of the library:
|
||||
- We follow PyTorch's API with methods like [`DiffusionPipeline.to`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.to) to let the user handle device management.
|
||||
- Raising concise error messages is preferred to silently correct erroneous input. Diffusers aims at teaching the user, rather than making the library as easy to use as possible.
|
||||
- Complex model vs. scheduler logic is exposed instead of magically handled inside. Schedulers/Samplers are separated from diffusion models with minimal dependencies on each other. This forces the user to write the unrolled denoising loop. However, the separation allows for easier debugging and gives the user more control over adapting the denoising process or switching out diffusion models or schedulers.
|
||||
- Separately trained components of the diffusion pipeline, *e.g.* the text encoder, the unet, and the variational autoencoder, each have their own model class. This forces the user to handle the interaction between the different model components, and the serialization format separates the model components into different files. However, this allows for easier debugging and customization. Dreambooth or textual inversion training
|
||||
is very simple thanks to diffusers' ability to separate single components of the diffusion pipeline.
|
||||
|
||||
## Tweakable, contributor-friendly over abstraction
|
||||
|
||||
For large parts of the library, Diffusers adopts an important design principle of the [Transformers library](https://github.com/huggingface/transformers), which is to prefer copy-pasted code over hasty abstractions. This design principle is very opinionated and stands in stark contrast to popular design principles such as [Don't repeat yourself (DRY)](https://en.wikipedia.org/wiki/Don%27t_repeat_yourself).
|
||||
In short, just like Transformers does for modeling files, diffusers prefers to keep an extremely low level of abstraction and very self-contained code for pipelines and schedulers.
|
||||
Functions, long code blocks, and even classes can be copied across multiple files which at first can look like a bad, sloppy design choice that makes the library unmaintainable.
|
||||
**However**, this design has proven to be extremely successful for Transformers and makes a lot of sense for community-driven, open-source machine learning libraries because:
|
||||
- Machine Learning is an extremely fast-moving field in which paradigms, model architectures, and algorithms are changing rapidly, which therefore makes it very difficult to define long-lasting code abstractions.
|
||||
- Machine Learning practitioners like to be able to quickly tweak existing code for ideation and research and therefore prefer self-contained code over one that contains many abstractions.
|
||||
- Open-source libraries rely on community contributions and therefore must build a library that is easy to contribute to. The more abstract the code, the more dependencies, the harder to read, and the harder to contribute to. Contributors simply stop contributing to very abstract libraries out of fear of breaking vital functionality. If contributing to a library cannot break other fundamental code, not only is it more inviting for potential new contributors, but it is also easier to review and contribute to multiple parts in parallel.
|
||||
|
||||
At Hugging Face, we call this design the **single-file policy** which means that almost all of the code of a certain class should be written in a single, self-contained file. To read more about the philosophy, you can have a look
|
||||
at [this blog post](https://huggingface.co/blog/transformers-design-philosophy).
|
||||
|
||||
In diffusers, we follow this philosophy for both pipelines and schedulers, but only partly for diffusion models. The reason we don't follow this design fully for diffusion models is because almost all diffusion pipelines, such
|
||||
as [DDPM](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/ddpm), [Stable Diffusion](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/stable_diffusion/overview#stable-diffusion-pipelines), [UnCLIP (Dalle-2)](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/unclip#overview) and [Imagen](https://imagen.research.google/) all rely on the same diffusion model, the [UNet](https://huggingface.co/docs/diffusers/api/models#diffusers.UNet2DConditionModel).
|
||||
|
||||
Great, now you should have generally understood why 🧨 Diffusers is designed the way it is 🤗.
|
||||
We try to apply these design principles consistently across the library. Nevertheless, there are some minor exceptions to the philosophy or some unlucky design choices. If you have feedback regarding the design, we would ❤️ to hear it [directly on GitHub](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feedback.md&title=).
|
||||
|
||||
## Design Philosophy in Details
|
||||
|
||||
Now, let's look a bit into the nitty-gritty details of the design philosophy. Diffusers essentially consist of three major classes, [pipelines](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines), [models](https://github.com/huggingface/diffusers/tree/main/src/diffusers/models), and [schedulers](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers).
|
||||
Let's walk through more in-detail design decisions for each class.
|
||||
|
||||
### Pipelines
|
||||
|
||||
Pipelines are designed to be easy to use (therefore do not follow [*Simple over easy*](#simple-over-easy) 100%)), are not feature complete, and should loosely be seen as examples of how to use [models](#models) and [schedulers](#schedulers) for inference.
|
||||
|
||||
The following design principles are followed:
|
||||
- Pipelines follow the single-file policy. All pipelines can be found in individual directories under src/diffusers/pipelines. One pipeline folder corresponds to one diffusion paper/project/release. Multiple pipeline files can be gathered in one pipeline folder, as it’s done for [`src/diffusers/pipelines/stable-diffusion`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines/stable_diffusion). If pipelines share similar functionality, one can make use of the [#Copied from mechanism](https://github.com/huggingface/diffusers/blob/125d783076e5bd9785beb05367a2d2566843a271/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_img2img.py#L251).
|
||||
- Pipelines all inherit from [`DiffusionPipeline`]
|
||||
- Every pipeline consists of different model and scheduler components, that are documented in the [`model_index.json` file](https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/model_index.json), are accessible under the same name as attributes of the pipeline and can be shared between pipelines with [`DiffusionPipeline.components`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.components) function.
|
||||
- Every pipeline should be loadable via the [`DiffusionPipeline.from_pretrained`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.from_pretrained) function.
|
||||
- Pipelines should be used **only** for inference.
|
||||
- Pipelines should be very readable, self-explanatory, and easy to tweak.
|
||||
- Pipelines should be designed to build on top of each other and be easy to integrate into higher-level APIs.
|
||||
- Pipelines are **not** intended to be feature-complete user interfaces. For future complete user interfaces one should rather have a look at [InvokeAI](https://github.com/invoke-ai/InvokeAI), [Diffuzers](https://github.com/abhishekkrthakur/diffuzers), and [lama-cleaner](https://github.com/Sanster/lama-cleaner)
|
||||
- Every pipeline should have one and only one way to run it via a `__call__` method. The naming of the `__call__` arguments should be shared across all pipelines.
|
||||
- Pipelines should be named after the task they are intended to solve.
|
||||
- In almost all cases, novel diffusion pipelines shall be implemented in a new pipeline folder/file.
|
||||
|
||||
### Models
|
||||
|
||||
Models are designed as configurable toolboxes that are natural extensions of [PyTorch's Module class](https://pytorch.org/docs/stable/generated/torch.nn.Module.html). They only partly follow the **single-file policy**.
|
||||
|
||||
The following design principles are followed:
|
||||
- Models correspond to **a type of model architecture**. *E.g.* the [`UNet2DConditionModel`] class is used for all UNet variations that expect 2D image inputs and are conditioned on some context.
|
||||
- All models can be found in [`src/diffusers/models`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/models) and every model architecture shall be defined in its file, e.g. [`unet_2d_condition.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_condition.py), [`transformer_2d.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/transformer_2d.py), etc...
|
||||
- Models **do not** follow the single-file policy and should make use of smaller model building blocks, such as [`attention.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention.py), [`resnet.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/resnet.py), [`embeddings.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/embeddings.py), etc... **Note**: This is in stark contrast to Transformers' modeling files and shows that models do not really follow the single-file policy.
|
||||
- Models intend to expose complexity, just like PyTorch's module does, and give clear error messages.
|
||||
- Models all inherit from `ModelMixin` and `ConfigMixin`.
|
||||
- Models can be optimized for performance when it doesn’t demand major code changes, keeps backward compatibility, and gives significant memory or compute gain.
|
||||
- Models should by default have the highest precision and lowest performance setting.
|
||||
- To integrate new model checkpoints whose general architecture can be classified as an architecture that already exists in Diffusers, the existing model architecture shall be adapted to make it work with the new checkpoint. One should only create a new file if the model architecture is fundamentally different.
|
||||
- Models should be designed to be easily extendable to future changes. This can be achieved by limiting public function arguments, configuration arguments, and "foreseeing" future changes, *e.g.* it is usually better to add `string` "...type" arguments that can easily be extended to new future types instead of boolean `is_..._type` arguments. Only the minimum amount of changes shall be made to existing architectures to make a new model checkpoint work.
|
||||
- The model design is a difficult trade-off between keeping code readable and concise and supporting many model checkpoints. For most parts of the modeling code, classes shall be adapted for new model checkpoints, while there are some exceptions where it is preferred to add new classes to make sure the code is kept concise and
|
||||
readable longterm, such as [UNet blocks](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_blocks.py) and [Attention processors](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/cross_attention.py).
|
||||
|
||||
### Schedulers
|
||||
|
||||
Schedulers are responsible to guide the denoising process for inference as well as to define a noise schedule for training. They are designed as individual classes with loadable configuration files and strongly follow the **single-file policy**.
|
||||
|
||||
The following design principles are followed:
|
||||
- All schedulers are found in [`src/diffusers/schedulers`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers).
|
||||
- Schedulers are **not** allowed to import from large utils files and shall be kept very self-contained.
|
||||
- One scheduler python file corresponds to one scheduler algorithm (as might be defined in a paper).
|
||||
- If schedulers share similar functionalities, we can make use of the `#Copied from` mechanism.
|
||||
- Schedulers all inherit from `SchedulerMixin` and `ConfigMixin`.
|
||||
- Schedulers can be easily swapped out with the [`ConfigMixin.from_config`](https://huggingface.co/docs/diffusers/main/en/api/configuration#diffusers.ConfigMixin.from_config) method as explained in detail [here](./using-diffusers/schedulers.mdx).
|
||||
- Every scheduler has to have a `set_num_inference_steps`, and a `step` function. `set_num_inference_steps(...)` has to be called before every denoising process, *i.e.* before `step(...)` is called.
|
||||
- Every scheduler exposes the timesteps to be "looped over" via a `timesteps` attribute, which is an array of timesteps the model will be called upon
|
||||
- The `step(...)` function takes a predicted model output and the "current" sample (x_t) and returns the "previous", slightly more denoised sample (x_t-1).
|
||||
- Given the complexity of diffusion schedulers, the `step` function does not expose all the complexity and can be a bit of a "black box".
|
||||
- In almost all cases, novel schedulers shall be implemented in a new scheduling file.
|
||||
Binary file not shown.
|
Before Width: | Height: | Size: 102 KiB |
@@ -1,75 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
<p align="center">
|
||||
<br>
|
||||
<img src="https://raw.githubusercontent.com/huggingface/diffusers/77aadfee6a891ab9fcfb780f87c693f7a5beeb8e/docs/source/imgs/diffusers_library.jpg" width="400"/>
|
||||
<br>
|
||||
</p>
|
||||
|
||||
# 🧨 Diffusers
|
||||
|
||||
🤗 Diffusers provides pretrained vision and audio diffusion models, and serves as a modular toolbox for inference and training.
|
||||
|
||||
More precisely, 🤗 Diffusers offers:
|
||||
|
||||
- State-of-the-art diffusion pipelines that can be run in inference with just a couple of lines of code (see [**Using Diffusers**](./using-diffusers/conditional_image_generation)) or have a look at [**Pipelines**](#pipelines) to get an overview of all supported pipelines and their corresponding papers.
|
||||
- Various noise schedulers that can be used interchangeably for the preferred speed vs. quality trade-off in inference. For more information see [**Schedulers**](./api/schedulers/overview).
|
||||
- Multiple types of models, such as UNet, can be used as building blocks in an end-to-end diffusion system. See [**Models**](./api/models) for more details
|
||||
- Training examples to show how to train the most popular diffusion model tasks. For more information see [**Training**](./training/overview).
|
||||
|
||||
## 🧨 Diffusers Pipelines
|
||||
|
||||
The following table summarizes all officially supported pipelines, their corresponding paper, and if
|
||||
available a colab notebook to directly try them out.
|
||||
|
||||
| Pipeline | Paper | Tasks | Colab
|
||||
|---|---|:---:|:---:|
|
||||
| [alt_diffusion](./api/pipelines/alt_diffusion) | [**AltDiffusion**](https://arxiv.org/abs/2211.06679) | Image-to-Image Text-Guided Generation |
|
||||
| [audio_diffusion](./api/pipelines/audio_diffusion) | [**Audio Diffusion**](https://github.com/teticio/audio-diffusion.git) | Unconditional Audio Generation | [](https://colab.research.google.com/github/teticio/audio-diffusion/blob/master/notebooks/audio_diffusion_pipeline.ipynb)
|
||||
| [cycle_diffusion](./api/pipelines/cycle_diffusion) | [**Cycle Diffusion**](https://arxiv.org/abs/2210.05559) | Image-to-Image Text-Guided Generation |
|
||||
| [dance_diffusion](./api/pipelines/dance_diffusion) | [**Dance Diffusion**](https://github.com/williamberman/diffusers.git) | Unconditional Audio Generation |
|
||||
| [ddpm](./api/pipelines/ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) | Unconditional Image Generation |
|
||||
| [ddim](./api/pipelines/ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) | Unconditional Image Generation |
|
||||
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation |
|
||||
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image |
|
||||
| [latent_diffusion_uncond](./api/pipelines/latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation |
|
||||
| [paint_by_example](./api/pipelines/paint_by_example) | [**Paint by Example: Exemplar-based Image Editing with Diffusion Models**](https://arxiv.org/abs/2211.13227) | Image-Guided Image Inpainting |
|
||||
| [pndm](./api/pipelines/pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation |
|
||||
| [score_sde_ve](./api/pipelines/score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
|
||||
| [score_sde_vp](./api/pipelines/score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
|
||||
| [semantic_stable_diffusion](./api/pipelines/semantic_stable_diffusion) | [**Semantic Guidance**](https://arxiv.org/abs/2301.12247) | Text-Guided Generation | [](https://colab.research.google.com/github/ml-research/semantic-image-editing/blob/main/examples/SemanticGuidance.ipynb)
|
||||
| [stable_diffusion_text2img](./api/pipelines/stable_diffusion/text2img) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
|
||||
| [stable_diffusion_img2img](./api/pipelines/stable_diffusion/img2img) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
|
||||
| [stable_diffusion_inpaint](./api/pipelines/stable_diffusion/inpaint) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
|
||||
| [stable_diffusion_panorama](./api/pipelines/stable_diffusion/panorama) | [**MultiDiffusion**](https://multidiffusion.github.io/) | Text-to-Panorama Generation |
|
||||
| [stable_diffusion_pix2pix](./api/pipelines/stable_diffusion/pix2pix) | [**InstructPix2Pix**](https://github.com/timothybrooks/instruct-pix2pix) | Text-Guided Image Editing|
|
||||
| [stable_diffusion_pix2pix_zero](./api/pipelines/stable_diffusion/pix2pix_zero) | [**Zero-shot Image-to-Image Translation**](https://pix2pixzero.github.io/) | Text-Guided Image Editing |
|
||||
| [stable_diffusion_attend_and_excite](./api/pipelines/stable_diffusion/attend_and_excite) | [**Attend and Excite for Stable Diffusion**](https://attendandexcite.github.io/Attend-and-Excite/) | Text-to-Image Generation |
|
||||
| [stable_diffusion_self_attention_guidance](./api/pipelines/stable_diffusion/self_attention_guidance) | [**Self-Attention Guidance**](https://ku-cvlab.github.io/Self-Attention-Guidance) | Text-to-Image Generation |
|
||||
| [stable_diffusion_image_variation](./stable_diffusion/image_variation) | [**Stable Diffusion Image Variations**](https://github.com/LambdaLabsML/lambda-diffusers#stable-diffusion-image-variations) | Image-to-Image Generation |
|
||||
| [stable_diffusion_latent_upscale](./stable_diffusion/latent_upscale) | [**Stable Diffusion Latent Upscaler**](https://twitter.com/StabilityAI/status/1590531958815064065) | Text-Guided Super Resolution Image-to-Image |
|
||||
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-to-Image Generation |
|
||||
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting |
|
||||
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Depth-Conditional Stable Diffusion**](https://github.com/Stability-AI/stablediffusion#depth-conditional-stable-diffusion) | Depth-to-Image Generation |
|
||||
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
|
||||
| [stable_diffusion_safe](./api/pipelines/stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
|
||||
| [stable_unclip](./stable_unclip) | **Stable unCLIP** | Text-to-Image Generation |
|
||||
| [stable_unclip](./stable_unclip) | **Stable unCLIP** | Image-to-Image Text-Guided Generation |
|
||||
| [stochastic_karras_ve](./api/pipelines/stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
|
||||
| [unclip](./api/pipelines/unclip) | [Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) | Text-to-Image Generation |
|
||||
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
|
||||
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation |
|
||||
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |
|
||||
| [vq_diffusion](./api/pipelines/vq_diffusion) | [Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation |
|
||||
|
||||
**Note**: Pipelines are simple examples of how to play around with the diffusion systems as described in the corresponding papers.
|
||||
@@ -1,70 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# How to use Stable Diffusion on Habana Gaudi
|
||||
|
||||
🤗 Diffusers is compatible with Habana Gaudi through 🤗 [Optimum Habana](https://huggingface.co/docs/optimum/habana/usage_guides/stable_diffusion).
|
||||
|
||||
## Requirements
|
||||
|
||||
- Optimum Habana 1.3 or later, [here](https://huggingface.co/docs/optimum/habana/installation) is how to install it.
|
||||
- SynapseAI 1.7.
|
||||
|
||||
|
||||
## Inference Pipeline
|
||||
|
||||
To generate images with Stable Diffusion 1 and 2 on Gaudi, you need to instantiate two instances:
|
||||
- A pipeline with [`GaudiStableDiffusionPipeline`](https://huggingface.co/docs/optimum/habana/package_reference/stable_diffusion_pipeline). This pipeline supports *text-to-image generation*.
|
||||
- A scheduler with [`GaudiDDIMScheduler`](https://huggingface.co/docs/optimum/habana/package_reference/stable_diffusion_pipeline#optimum.habana.diffusers.GaudiDDIMScheduler). This scheduler has been optimized for Habana Gaudi.
|
||||
|
||||
When initializing the pipeline, you have to specify `use_habana=True` to deploy it on HPUs.
|
||||
Furthermore, in order to get the fastest possible generations you should enable **HPU graphs** with `use_hpu_graphs=True`.
|
||||
Finally, you will need to specify a [Gaudi configuration](https://huggingface.co/docs/optimum/habana/package_reference/gaudi_config) which can be downloaded from the [Hugging Face Hub](https://huggingface.co/Habana).
|
||||
|
||||
```python
|
||||
from optimum.habana import GaudiConfig
|
||||
from optimum.habana.diffusers import GaudiDDIMScheduler, GaudiStableDiffusionPipeline
|
||||
|
||||
model_name = "stabilityai/stable-diffusion-2-base"
|
||||
scheduler = GaudiDDIMScheduler.from_pretrained(model_name, subfolder="scheduler")
|
||||
pipeline = GaudiStableDiffusionPipeline.from_pretrained(
|
||||
model_name,
|
||||
scheduler=scheduler,
|
||||
use_habana=True,
|
||||
use_hpu_graphs=True,
|
||||
gaudi_config="Habana/stable-diffusion",
|
||||
)
|
||||
```
|
||||
|
||||
You can then call the pipeline to generate images by batches from one or several prompts:
|
||||
```python
|
||||
outputs = pipeline(
|
||||
prompt=[
|
||||
"High quality photo of an astronaut riding a horse in space",
|
||||
"Face of a yellow cat, high resolution, sitting on a park bench",
|
||||
],
|
||||
num_images_per_prompt=10,
|
||||
batch_size=4,
|
||||
)
|
||||
```
|
||||
|
||||
For more information, check out Optimum Habana's [documentation](https://huggingface.co/docs/optimum/habana/usage_guides/stable_diffusion) and the [example](https://github.com/huggingface/optimum-habana/tree/main/examples/stable-diffusion) provided in the official Github repository.
|
||||
|
||||
|
||||
## Benchmark
|
||||
|
||||
Here are the latencies for Habana Gaudi 1 and Gaudi 2 with the [Habana/stable-diffusion](https://huggingface.co/Habana/stable-diffusion) Gaudi configuration (mixed precision bf16/fp32):
|
||||
|
||||
| | Latency | Batch size |
|
||||
| ------- |:-------:|:----------:|
|
||||
| Gaudi 1 | 4.37s | 4/8 |
|
||||
| Gaudi 2 | 1.19s | 4/8 |
|
||||
@@ -1,200 +0,0 @@
|
||||
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Torch2.0 support in Diffusers
|
||||
|
||||
Starting from version `0.13.0`, Diffusers supports the latest optimization from the upcoming [PyTorch 2.0](https://pytorch.org/get-started/pytorch-2.0/) release. These include:
|
||||
1. Support for native flash and memory-efficient attention without any extra dependencies.
|
||||
2. [torch.compile](https://pytorch.org/tutorials/intermediate/torch_compile_tutorial.html) support for compiling individual models for extra performance boost.
|
||||
|
||||
|
||||
## Installation
|
||||
To benefit from the native efficient attention and `torch.compile`, we will need to install the nightly version of PyTorch as the stable version is yet to be released. The first step is to install CUDA11.7 or CUDA11.8,
|
||||
as torch2.0 does not support the previous versions. Once CUDA is installed, torch nightly can be installed using:
|
||||
|
||||
```bash
|
||||
pip install --pre torch torchvision --index-url https://download.pytorch.org/whl/nightly/cu117
|
||||
```
|
||||
|
||||
## Using efficient attention and torch.compile.
|
||||
|
||||
|
||||
1. **Efficient Attention**
|
||||
|
||||
Efficient attention is implemented via the [`torch.nn.functional.scaled_dot_product_attention`](https://pytorch.org/docs/master/generated/torch.nn.functional.scaled_dot_product_attention) function, which automatically enables flash/memory efficient attention, depending on the input and the GPU type. This is the same as the `memory_efficient_attention` from [xFormers](https://github.com/facebookresearch/xformers) but built natively into PyTorch.
|
||||
|
||||
Efficient attention will be enabled by default in Diffusers if torch2.0 is installed and if `torch.nn.functional.scaled_dot_product_attention` is available. To use it, you can install torch2.0 as suggested above and use the pipeline. For example:
|
||||
|
||||
```Python
|
||||
import torch
|
||||
from diffusers import StableDiffusionPipeline
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "a photo of an astronaut riding a horse on mars"
|
||||
image = pipe(prompt).images[0]
|
||||
```
|
||||
|
||||
If you want to enable it explicitly (which is not required), you can do so as shown below.
|
||||
|
||||
```Python
|
||||
import torch
|
||||
from diffusers import StableDiffusionPipeline
|
||||
from diffusers.models.cross_attention import AttnProcessor2_0
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16).to("cuda")
|
||||
pipe.unet.set_attn_processor(AttnProcessor2_0())
|
||||
|
||||
prompt = "a photo of an astronaut riding a horse on mars"
|
||||
image = pipe(prompt).images[0]
|
||||
```
|
||||
|
||||
This should be as fast and memory efficient as `xFormers`.
|
||||
|
||||
|
||||
2. **torch.compile**
|
||||
|
||||
To get an additional speedup, we can use the new `torch.compile` feature. To do so, we wrap our `unet` with `torch.compile`. For more information and different options, refer to the
|
||||
[torch compile docs](https://pytorch.org/tutorials/intermediate/torch_compile_tutorial.html).
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import StableDiffusionPipeline
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16).to(
|
||||
"cuda"
|
||||
)
|
||||
pipe.unet = torch.compile(pipe.unet)
|
||||
|
||||
batch_size = 10
|
||||
prompt = "A photo of an astronaut riding a horse on marse."
|
||||
images = pipe(prompt, num_inference_steps=steps, num_images_per_prompt=batch_size).images
|
||||
```
|
||||
|
||||
Depending on the type of GPU it can give between 2-9% speed-up over efficient attention. But note that as of now the speed-up is mostly noticeable on the more recent GPU architectures, such as in the A100.
|
||||
|
||||
Note that compilation will also take some time to complete, so it is best suited for situations where you need to prepare your pipeline once and then perform the same type of inference operations multiple times.
|
||||
|
||||
|
||||
## Benchmark
|
||||
|
||||
We conducted a simple benchmark on different GPUs to compare vanilla attention, xFormers, `torch.nn.functional.scaled_dot_product_attention` and `torch.compile+torch.nn.functional.scaled_dot_product_attention`.
|
||||
For the benchmark we used the the [stable-diffusion-v1-4](https://huggingface.co/CompVis/stable-diffusion-v1-4) model with 50 steps. `xFormers` benchmark is done using the `torch==1.13.1` version. The table below summarizes the result that we got.
|
||||
The `Speed over xformers` columns denotes the speed-up gained over `xFormers` using the `torch.compile+torch.nn.functional.scaled_dot_product_attention`.
|
||||
|
||||
|
||||
### FP16 benchmark
|
||||
|
||||
The table below shows the benchmark results for inference using `fp16`. As we can see, `torch.nn.functional.scaled_dot_product_attention` is as fast as `xFormers` (sometimes slightly faster/slower) on all the GPUs we tested.
|
||||
And using `torch.compile` gives further speed-up up to 10% over `xFormers`, but it's mostly noticeable on the A100 GPU.
|
||||
|
||||
___The time reported is in seconds.___
|
||||
|
||||
| GPU | Batch Size | Vanilla Attention | xFormers | PyTorch2.0 SDPA | SDPA + torch.compile | Speed over xformers (%) |
|
||||
| --- | --- | --- | --- | --- | --- | --- |
|
||||
| A100 | 10 | 12.02 | 8.7 | 8.79 | 7.89 | 9.31 |
|
||||
| A100 | 16 | 18.95 | 13.57 | 13.67 | 12.25 | 9.73 |
|
||||
| A100 | 32 (1) | OOM | 26.56 | 26.68 | 24.08 | 9.34 |
|
||||
| A100 | 64(2) | | 52.51 | 53.03 | 47.81 | 8.95 |
|
||||
| | | | | | | |
|
||||
| A10 | 4 | 13.94 | 9.81 | 10.01 | 9.35 | 4.69 |
|
||||
| A10 | 8 | 27.09 | 19 | 19.53 | 18.33 | 3.53 |
|
||||
| A10 | 10 | 33.69 | 23.53 | 24.19 | 22.52 | 4.29 |
|
||||
| A10 | 16 | OOM | 37.55 | 38.31 | 36.81 | 1.97 |
|
||||
| A10 | 32 (1) | | 77.19 | 78.43 | 76.64 | 0.71 |
|
||||
| A10 | 64 (1) | | 173.59 | 158.99 | 155.14 | 10.63 |
|
||||
| | | | | | | |
|
||||
| T4 | 4 | 38.81 | 30.09 | 29.74 | 27.55 | 8.44 |
|
||||
| T4 | 8 | OOM | 55.71 | 55.99 | 53.85 | 3.34 |
|
||||
| T4 | 10 | OOM | 68.96 | 69.86 | 65.35 | 5.23 |
|
||||
| T4 | 16 | OOM | 111.47 | 113.26 | 106.93 | 4.07 |
|
||||
| | | | | | | |
|
||||
| V100 | 4 | 9.84 | 8.16 | 8.09 | 7.65 | 6.25 |
|
||||
| V100 | 8 | OOM | 15.62 | 15.44 | 14.59 | 6.59 |
|
||||
| V100 | 10 | OOM | 19.52 | 19.28 | 18.18 | 6.86 |
|
||||
| V100 | 16 | OOM | 30.29 | 29.84 | 28.22 | 6.83 |
|
||||
| | | | | | | |
|
||||
| 3090 | 4 | 10.04 | 7.82 | 7.89 | 7.47 | 4.48 |
|
||||
| 3090 | 8 | 19.27 | 14.97 | 15.04 | 14.22 | 5.01 |
|
||||
| 3090 | 10| 24.08 | 18.7 | 18.7 | 17.69 | 5.40 |
|
||||
| 3090 | 16 | OOM | 29.06 | 29.06 | 28.2 | 2.96 |
|
||||
| 3090 | 32 (1) | | 58.05 | 58 | 54.88 | 5.46 |
|
||||
| 3090 | 64 (1) | | 126.54 | 126.03 | 117.33 | 7.28 |
|
||||
| | | | | | | |
|
||||
| 3090 Ti | 4 | 9.07 | 7.14 | 7.15 | 6.81 | 4.62 |
|
||||
| 3090 Ti | 8 | 17.51 | 13.65 | 13.72 | 12.99 | 4.84 |
|
||||
| 3090 Ti | 10 (2) | 21.79 | 16.85 | 16.93 | 16.02 | 4.93 |
|
||||
| 3090 Ti | 16 | OOM | 26.1 | 26.28 | 25.46 | 2.45 |
|
||||
| 3090 Ti | 32 (1) | | 51.78 | 52.04 | 49.15 | 5.08 |
|
||||
| 3090 Ti | 64 (1) | | 112.02 | 112.33 | 103.91 | 7.24 |
|
||||
|
||||
|
||||
|
||||
### FP32 benchmark
|
||||
|
||||
The table below shows the benchmark results for inference using `fp32`. As we can see, `torch.nn.functional.scaled_dot_product_attention` is as fast as `xFormers` (sometimes slightly faster/slower) on all the GPUs we tested.
|
||||
Using `torch.compile` with efficient attention gives up to 18% performance improvement over `xFormers` in Ampere cards, and up to 20% over vanilla attention.
|
||||
|
||||
| GPU | Batch Size | Vanilla Attention | xFormers | PyTorch2.0 SDPA | SDPA + torch.compile | Speed over xformers (%) | Speed over vanilla (%) |
|
||||
| --- | --- | --- | --- | --- | --- | --- | --- |
|
||||
| A100 | 4 | 16.56 | 12.42 | 12.2 | 11.84 | 4.67 | 28.50 |
|
||||
| A100 | 10 | OOM | 29.93 | 29.44 | 28.5 | 4.78 | |
|
||||
| A100 | 16 | | 47.08 | 46.27 | 44.8 | 4.84 | |
|
||||
| A100 | 32 | | 92.89 | 91.34 | 88.35 | 4.89 | |
|
||||
| A100 | 64 | | 185.3 | 182.71 | 176.48 | 4.76 | |
|
||||
| | | | | | | |
|
||||
| A10 | 1 | 10.59 | 8.81 | 7.51 | 7.35 | 16.57 | 30.59 |
|
||||
| A10 | 4 | 34.77 | 27.63 | 22.77 | 22.07 | 20.12 | 36.53 |
|
||||
| A10 | 8 | | 56.19 | 43.53 | 43.86 | 21.94 | |
|
||||
| A10 | 16 | | 116.49 | 88.56 | 86.64 | 25.62 | |
|
||||
| A10 | 32 | | 221.95 | 175.74 | 168.18 | 24.23 | |
|
||||
| A10 | 48 | | 333.23 | 264.84 | | 20.52 | |
|
||||
| | | | | | | |
|
||||
| T4 | 1 | 28.2 | 24.49 | 23.93 | 23.56 | 3.80 | 16.45 |
|
||||
| T4 | 2 | 52.77 | 45.7 | 45.88 | 45.06 | 1.40 | 14.61 |
|
||||
| T4 | 4 | OOM | 85.72 | 85.78 | 84.48 | 1.45 | |
|
||||
| T4 | 8 | | 149.64 | 150.75 | 148.4 | 0.83 | |
|
||||
| | | | | | | |
|
||||
| V100 | 1 | 7.4 | 6.84 | 6.8 | 6.66 | 2.63 | 10.00 |
|
||||
| V100 | 2 | 13.85 | 12.81 | 12.66 | 12.35 | 3.59 | 10.83 |
|
||||
| V100 | 4 | OOM | 25.73 | 25.31 | 24.78 | 3.69 | |
|
||||
| V100 | 8 | | 43.95 | 43.37 | 42.25 | 3.87 | |
|
||||
| V100 | 16 | | 84.99 | 84.73 | 82.55 | 2.87 | |
|
||||
| | | | | | | |
|
||||
| 3090 | 1 | 7.09 | 6.78 | 6.11 | 6.03 | 11.06 | 14.95 |
|
||||
| 3090 | 4 | 22.69 | 21.45 | 18.67 | 18.09 | 15.66 | 20.27 |
|
||||
| 3090 | 8 (2) | | 42.59 | 36.75 | 35.59 | 16.44 | |
|
||||
| 3090 | 16 | | 85.35 | 72.37 | 70.25 | 17.69 | |
|
||||
| 3090 | 32 (1) | | 162.05 | 138.99 | 134.53 | 16.98 | |
|
||||
| 3090 | 48 | | 241.91 | 207.75 | | 14.12 | |
|
||||
| | | | | | | |
|
||||
| 3090 Ti | 1 | 6.45 | 6.19 | 5.64 | 5.49 | 11.31 | 14.88 |
|
||||
| 3090 Ti | 4 | 20.32 | 19.31 | 16.9 | 16.37 | 15.23 | 19.44 |
|
||||
| 3090 Ti | 8 (2) | | 37.93 | 33.05 | 31.99 | 15.66 | |
|
||||
| 3090 Ti | 16 | | 75.37 | 65.25 | 64.32 | 14.66 | |
|
||||
| 3090 Ti | 32 (1) | | 142.55 | 124.44 | 120.74 | 15.30 | |
|
||||
| 3090 Ti | 48 | | 213.19 | 186.55 | | 12.50 | |
|
||||
| | | | | | | |
|
||||
| 4090 | 1 | 5.54 | 4.99 | 4.51 | | | |
|
||||
| 4090 | 4 | 13.67 | 11.4 | 10.3 | | | |
|
||||
| 4090 | 8 (2) | | 19.79 | 17.13 | | | |
|
||||
| 4090 | 16 | | 38.62 | 33.14 | | | |
|
||||
| 4090 | 32 (1) | | 76.57 | 65.96 | | | |
|
||||
| 4090 | 48 | | 114.44 | 98.78 | | | |
|
||||
|
||||
|
||||
|
||||
(1) Batch Size >= 32 requires enable_vae_slicing() because of https://github.com/pytorch/pytorch/issues/81665
|
||||
This is required for PyTorch 1.13.1, and also for PyTorch 2.0 and batch size of 64
|
||||
|
||||
For more details about how this benchmark was run, please refer to [this PR](https://github.com/huggingface/diffusers/pull/2303).
|
||||
@@ -1,35 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Installing xFormers
|
||||
|
||||
We recommend the use of [xFormers](https://github.com/facebookresearch/xformers) for both inference and training. In our tests, the optimizations performed in the attention blocks allow for both faster speed and reduced memory consumption.
|
||||
|
||||
Starting from version `0.0.16` of xFormers, released on January 2023, installation can be easily performed using pre-built pip wheels:
|
||||
|
||||
```bash
|
||||
pip install xformers
|
||||
```
|
||||
|
||||
<Tip>
|
||||
|
||||
The xFormers PIP package requires the latest version of PyTorch (1.13.1 as of xFormers 0.0.16). If you need to use a previous version of PyTorch, then we recommend you install xFormers from source using [the project instructions](https://github.com/facebookresearch/xformers#installing-xformers).
|
||||
|
||||
</Tip>
|
||||
|
||||
After xFormers is installed, you can use `enable_xformers_memory_efficient_attention()` for faster inference and reduced memory consumption, as discussed [here](fp16#memory-efficient-attention).
|
||||
|
||||
<Tip warning={true}>
|
||||
|
||||
According to [this issue](https://github.com/huggingface/diffusers/issues/2234#issuecomment-1416931212), xFormers `v0.0.16` cannot be used for training (fine-tune or Dreambooth) in some GPUs. If you observe that problem, please install a development version as indicated in that comment.
|
||||
|
||||
</Tip>
|
||||
@@ -1,333 +0,0 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# The Stable Diffusion Guide 🎨
|
||||
<a target="_blank" href="https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/sd_101_guide.ipynb">
|
||||
<img src="https://colab.research.google.com/assets/colab-badge.svg" alt="Open In Colab"/>
|
||||
</a>
|
||||
|
||||
## Intro
|
||||
|
||||
Stable Diffusion is a [Latent Diffusion model](https://github.com/CompVis/latent-diffusion) developed by researchers from the Machine Vision and Learning group at LMU Munich, *a.k.a* CompVis.
|
||||
Model checkpoints were publicly released at the end of August 2022 by a collaboration of Stability AI, CompVis, and Runway with support from EleutherAI and LAION. For more information, you can check out [the official blog post](https://stability.ai/blog/stable-diffusion-public-release).
|
||||
|
||||
Since its public release the community has done an incredible job at working together to make the stable diffusion checkpoints **faster**, **more memory efficient**, and **more performant**.
|
||||
|
||||
🧨 Diffusers offers a simple API to run stable diffusion with all memory, computing, and quality improvements.
|
||||
|
||||
This notebook walks you through the improvements one-by-one so you can best leverage [`StableDiffusionPipeline`] for **inference**.
|
||||
|
||||
## Prompt Engineering 🎨
|
||||
|
||||
When running *Stable Diffusion* in inference, we usually want to generate a certain type, or style of image and then improve upon it. Improving upon a previously generated image means running inference over and over again with a different prompt and potentially a different seed until we are happy with our generation.
|
||||
|
||||
So to begin with, it is most important to speed up stable diffusion as much as possible to generate as many pictures as possible in a given amount of time.
|
||||
|
||||
This can be done by both improving the **computational efficiency** (speed) and the **memory efficiency** (GPU RAM).
|
||||
|
||||
Let's start by looking into computational efficiency first.
|
||||
|
||||
Throughout the notebook, we will focus on [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5):
|
||||
|
||||
``` python
|
||||
model_id = "runwayml/stable-diffusion-v1-5"
|
||||
```
|
||||
|
||||
Let's load the pipeline.
|
||||
|
||||
## Speed Optimization
|
||||
|
||||
``` python
|
||||
from diffusers import StableDiffusionPipeline
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained(model_id)
|
||||
```
|
||||
|
||||
We aim at generating a beautiful photograph of an *old warrior chief* and will later try to find the best prompt to generate such a photograph. For now, let's keep the prompt simple:
|
||||
|
||||
``` python
|
||||
prompt = "portrait photo of a old warrior chief"
|
||||
```
|
||||
|
||||
To begin with, we should make sure we run inference on GPU, so let's move the pipeline to GPU, just like you would with any PyTorch module.
|
||||
|
||||
``` python
|
||||
pipe = pipe.to("cuda")
|
||||
```
|
||||
|
||||
To generate an image, you should use the [~`StableDiffusionPipeline.__call__`] method.
|
||||
|
||||
To make sure we can reproduce more or less the same image in every call, let's make use of the generator. See the documentation on reproducibility [here](./conceptual/reproducibility) for more information.
|
||||
|
||||
``` python
|
||||
generator = torch.Generator("cuda").manual_seed(0)
|
||||
```
|
||||
|
||||
Now, let's take a spin on it.
|
||||
|
||||
``` python
|
||||
image = pipe(prompt, generator=generator).images[0]
|
||||
image
|
||||
```
|
||||
|
||||

|
||||
|
||||
Cool, this now took roughly 30 seconds on a T4 GPU (you might see faster inference if your allocated GPU is better than a T4).
|
||||
|
||||
The default run we did above used full float32 precision and ran the default number of inference steps (50). The easiest speed-ups come from switching to float16 (or half) precision and simply running fewer inference steps. Let's load the model now in float16 instead.
|
||||
|
||||
``` python
|
||||
import torch
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
|
||||
pipe = pipe.to("cuda")
|
||||
```
|
||||
|
||||
And we can again call the pipeline to generate an image.
|
||||
|
||||
``` python
|
||||
generator = torch.Generator("cuda").manual_seed(0)
|
||||
|
||||
image = pipe(prompt, generator=generator).images[0]
|
||||
image
|
||||
```
|
||||

|
||||
|
||||
Cool, this is almost three times as fast for arguably the same image quality.
|
||||
|
||||
We strongly suggest always running your pipelines in float16 as so far we have very rarely seen degradations in quality because of it.
|
||||
|
||||
Next, let's see if we need to use 50 inference steps or whether we could use significantly fewer. The number of inference steps is associated with the denoising scheduler we use. Choosing a more efficient scheduler could help us decrease the number of steps.
|
||||
|
||||
Let's have a look at all the schedulers the stable diffusion pipeline is compatible with.
|
||||
|
||||
``` python
|
||||
pipe.scheduler.compatibles
|
||||
```
|
||||
|
||||
```
|
||||
[diffusers.schedulers.scheduling_dpmsolver_singlestep.DPMSolverSinglestepScheduler,
|
||||
diffusers.schedulers.scheduling_lms_discrete.LMSDiscreteScheduler,
|
||||
diffusers.schedulers.scheduling_heun_discrete.HeunDiscreteScheduler,
|
||||
diffusers.schedulers.scheduling_pndm.PNDMScheduler,
|
||||
diffusers.schedulers.scheduling_euler_discrete.EulerDiscreteScheduler,
|
||||
diffusers.schedulers.scheduling_euler_ancestral_discrete.EulerAncestralDiscreteScheduler,
|
||||
diffusers.schedulers.scheduling_dpmsolver_multistep.DPMSolverMultistepScheduler,
|
||||
diffusers.schedulers.scheduling_ddpm.DDPMScheduler,
|
||||
diffusers.schedulers.scheduling_ddim.DDIMScheduler]
|
||||
```
|
||||
|
||||
Cool, that's a lot of schedulers.
|
||||
|
||||
🧨 Diffusers is constantly adding a bunch of novel schedulers/samplers that can be used with Stable Diffusion. For more information, we recommend taking a look at the official documentation [here](https://huggingface.co/docs/diffusers/main/en/api/schedulers/overview).
|
||||
|
||||
Alright, right now Stable Diffusion is using the `PNDMScheduler` which usually requires around 50 inference steps. However, other schedulers such as `DPMSolverMultistepScheduler` or `DPMSolverSinglestepScheduler` seem to get away with just 20 to 25 inference steps. Let's try them out.
|
||||
|
||||
You can set a new scheduler by making use of the [from_config](https://huggingface.co/docs/diffusers/main/en/api/configuration#diffusers.ConfigMixin.from_config) function.
|
||||
|
||||
``` python
|
||||
from diffusers import DPMSolverMultistepScheduler
|
||||
|
||||
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
|
||||
```
|
||||
|
||||
Now, let's try to reduce the number of inference steps to just 20.
|
||||
|
||||
``` python
|
||||
generator = torch.Generator("cuda").manual_seed(0)
|
||||
|
||||
image = pipe(prompt, generator=generator, num_inference_steps=20).images[0]
|
||||
image
|
||||
```
|
||||
|
||||

|
||||
|
||||
The image now does look a little different, but it's arguably still of equally high quality. We now cut inference time to just 4 seconds though 😍.
|
||||
|
||||
## Memory Optimization
|
||||
|
||||
Less memory used in generation indirectly implies more speed, since we're often trying to maximize how many images we can generate per second. Usually, the more images per inference run, the more images per second too.
|
||||
|
||||
The easiest way to see how many images we can generate at once is to simply try it out, and see when we get a *"Out-of-memory (OOM)"* error.
|
||||
|
||||
We can run batched inference by simply passing a list of prompts and generators. Let's define a quick function that generates a batch for us.
|
||||
|
||||
``` python
|
||||
def get_inputs(batch_size=1):
|
||||
generator = [torch.Generator("cuda").manual_seed(i) for i in range(batch_size)]
|
||||
prompts = batch_size * [prompt]
|
||||
num_inference_steps = 20
|
||||
|
||||
return {"prompt": prompts, "generator": generator, "num_inference_steps": num_inference_steps}
|
||||
```
|
||||
This function returns a list of prompts and a list of generators, so we can reuse the generator that produced a result we like.
|
||||
|
||||
We also need a method that allows us to easily display a batch of images.
|
||||
|
||||
``` python
|
||||
from PIL import Image
|
||||
|
||||
def image_grid(imgs, rows=2, cols=2):
|
||||
w, h = imgs[0].size
|
||||
grid = Image.new('RGB', size=(cols*w, rows*h))
|
||||
|
||||
for i, img in enumerate(imgs):
|
||||
grid.paste(img, box=(i%cols*w, i//cols*h))
|
||||
return grid
|
||||
```
|
||||
|
||||
Cool, let's see how much memory we can use starting with `batch_size=4`.
|
||||
|
||||
``` python
|
||||
images = pipe(**get_inputs(batch_size=4)).images
|
||||
image_grid(images)
|
||||
```
|
||||
|
||||

|
||||
|
||||
Going over a batch_size of 4 will error out in this notebook (assuming we are running it on a T4 GPU). Also, we can see we only generate slightly more images per second (3.75s/image) compared to 4s/image previously.
|
||||
|
||||
However, the community has found some nice tricks to improve the memory constraints further. After stable diffusion was released, the community found improvements within days and shared them freely over GitHub - open-source at its finest! I believe the original idea came from [this](https://github.com/basujindal/stable-diffusion/pull/117) GitHub thread.
|
||||
|
||||
By far most of the memory is taken up by the cross-attention layers. Instead of running this operation in batch, one can run it sequentially to save a significant amount of memory.
|
||||
|
||||
It can easily be enabled by calling `enable_attention_slicing` as is documented [here](https://huggingface.co/docs/diffusers/main/en/api/pipelines/stable_diffusion/text2img#diffusers.StableDiffusionPipeline.enable_attention_slicing).
|
||||
|
||||
``` python
|
||||
pipe.enable_attention_slicing()
|
||||
```
|
||||
|
||||
Great, now that attention slicing is enabled, let's try to double the batch size again, going for `batch_size=8`.
|
||||
|
||||
``` python
|
||||
images = pipe(**get_inputs(batch_size=8)).images
|
||||
image_grid(images, rows=2, cols=4)
|
||||
```
|
||||
|
||||

|
||||
|
||||
Nice, it works. However, the speed gain is again not very big (it might however be much more significant on other GPUs).
|
||||
|
||||
We're at roughly 3.5 seconds per image 🔥 which is probably the fastest we can be with a simple T4 without sacrificing quality.
|
||||
|
||||
Next, let's look into how to improve the quality!
|
||||
|
||||
## Quality Improvements
|
||||
|
||||
Now that our image generation pipeline is blazing fast, let's try to get maximum image quality.
|
||||
|
||||
First of all, image quality is extremely subjective, so it's difficult to make general claims here.
|
||||
|
||||
The most obvious step to take to improve quality is to use *better checkpoints*. Since the release of Stable Diffusion, many improved versions have been released, which are summarized here:
|
||||
|
||||
- *Official Release - 22 Aug 2022*: [Stable-Diffusion 1.4](https://huggingface.co/CompVis/stable-diffusion-v1-4)
|
||||
- *20 October 2022*: [Stable-Diffusion 1.5](https://huggingface.co/runwayml/stable-diffusion-v1-5)
|
||||
- *24 Nov 2022*: [Stable-Diffusion 2.0](https://huggingface.co/stabilityai/stable-diffusion-2-0)
|
||||
- *7 Dec 2022*: [Stable-Diffusion 2.1](https://huggingface.co/stabilityai/stable-diffusion-2-1)
|
||||
|
||||
Newer versions don't necessarily mean better image quality with the same parameters. People mentioned that *2.0* is slightly worse than *1.5* for certain prompts, but given the right prompt engineering *2.0* and *2.1* seem to be better.
|
||||
|
||||
Overall, we strongly recommend just trying the models out and reading up on advice online (e.g. it has been shown that using negative prompts is very important for 2.0 and 2.1 to get the highest possible quality. See for example [this nice blog post](https://minimaxir.com/2022/11/stable-diffusion-negative-prompt/).
|
||||
|
||||
Additionally, the community has started fine-tuning many of the above versions on certain styles with some of them having an extremely high quality and gaining a lot of traction.
|
||||
|
||||
We recommend having a look at all [diffusers checkpoints sorted by downloads and trying out the different checkpoints](https://huggingface.co/models?library=diffusers).
|
||||
|
||||
For the following, we will stick to v1.5 for simplicity.
|
||||
|
||||
Next, we can also try to optimize single components of the pipeline, e.g. switching out the latent decoder. For more details on how the whole Stable Diffusion pipeline works, please have a look at [this blog post](https://huggingface.co/blog/stable_diffusion).
|
||||
|
||||
Let's load [stabilityai's newest auto-decoder](https://huggingface.co/stabilityai/stable-diffusion-2-1).
|
||||
|
||||
``` python
|
||||
from diffusers import AutoencoderKL
|
||||
|
||||
vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse", torch_dtype=torch.float16).to("cuda")
|
||||
```
|
||||
|
||||
Now we can set it to the vae of the pipeline to use it.
|
||||
|
||||
``` python
|
||||
pipe.vae = vae
|
||||
```
|
||||
|
||||
Let's run the same prompt as before to compare quality.
|
||||
|
||||
``` python
|
||||
images = pipe(**get_inputs(batch_size=8)).images
|
||||
image_grid(images, rows=2, cols=4)
|
||||
```
|
||||
|
||||

|
||||
|
||||
Seems like the difference is only very minor, but the new generations are arguably a bit *sharper*.
|
||||
|
||||
Cool, finally, let's look a bit into prompt engineering.
|
||||
|
||||
Our goal was to generate a photo of an old warrior chief. Let's now try to bring a bit more color into the photos and make the look more impressive.
|
||||
|
||||
Originally our prompt was "*portrait photo of an old warrior chief*".
|
||||
|
||||
To improve the prompt, it often helps to add cues that could have been used online to save high-quality photos, as well as add more details.
|
||||
Essentially, when doing prompt engineering, one has to think:
|
||||
|
||||
- How was the photo or similar photos of the one I want probably stored on the internet?
|
||||
- What additional detail can I give that steers the models into the style that I want?
|
||||
|
||||
Cool, let's add more details.
|
||||
|
||||
``` python
|
||||
prompt += ", tribal panther make up, blue on red, side profile, looking away, serious eyes"
|
||||
```
|
||||
|
||||
and let's also add some cues that usually help to generate higher quality images.
|
||||
|
||||
``` python
|
||||
prompt += " 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta"
|
||||
prompt
|
||||
```
|
||||
|
||||
Cool, let's now try this prompt.
|
||||
|
||||
``` python
|
||||
images = pipe(**get_inputs(batch_size=8)).images
|
||||
image_grid(images, rows=2, cols=4)
|
||||
```
|
||||
|
||||

|
||||
|
||||
Pretty impressive! We got some very high-quality image generations there. The 2nd image is my personal favorite, so I'll re-use this seed and see whether I can tweak the prompts slightly by using "oldest warrior", "old", "", and "young" instead of "old".
|
||||
|
||||
``` python
|
||||
prompts = [
|
||||
"portrait photo of the oldest warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
|
||||
"portrait photo of a old warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
|
||||
"portrait photo of a warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
|
||||
"portrait photo of a young warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
|
||||
]
|
||||
|
||||
generator = [torch.Generator("cuda").manual_seed(1) for _ in range(len(prompts))] # 1 because we want the 2nd image
|
||||
|
||||
images = pipe(prompt=prompts, generator=generator, num_inference_steps=25).images
|
||||
image_grid(images)
|
||||
```
|
||||
|
||||

|
||||
|
||||
The first picture looks nice! The eye movement slightly changed and looks nice. This finished up our 101-guide on how to use Stable Diffusion 🤗.
|
||||
|
||||
For more information on optimization or other guides, I recommend taking a look at the following:
|
||||
|
||||
- [Blog post about Stable Diffusion](https://huggingface.co/blog/stable_diffusion): In-detail blog post explaining Stable Diffusion.
|
||||
- [FlashAttention](https://huggingface.co/docs/diffusers/optimization/xformers): XFormers flash attention can optimize your model even further with more speed and memory improvements.
|
||||
- [Dreambooth](https://huggingface.co/docs/diffusers/training/dreambooth) - Quickly customize the model by fine-tuning it.
|
||||
- [General info on Stable Diffusion](https://huggingface.co/docs/diffusers/main/en/api/pipelines/stable_diffusion/overview) - Info on other tasks that are powered by Stable Diffusion.
|
||||
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Reference in New Issue
Block a user