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38
.github/ISSUE_TEMPLATE/remote-vae-pilot-feedback.yml
vendored
Normal file
38
.github/ISSUE_TEMPLATE/remote-vae-pilot-feedback.yml
vendored
Normal file
@@ -0,0 +1,38 @@
|
||||
name: "\U0001F31F Remote VAE"
|
||||
description: Feedback for remote VAE pilot
|
||||
labels: [ "Remote VAE" ]
|
||||
|
||||
body:
|
||||
- type: textarea
|
||||
id: positive
|
||||
validations:
|
||||
required: true
|
||||
attributes:
|
||||
label: Did you like the remote VAE solution?
|
||||
description: |
|
||||
If you liked it, we would appreciate it if you could elaborate what you liked.
|
||||
|
||||
- type: textarea
|
||||
id: feedback
|
||||
validations:
|
||||
required: true
|
||||
attributes:
|
||||
label: What can be improved about the current solution?
|
||||
description: |
|
||||
Let us know the things you would like to see improved. Note that we will work optimizing the solution once the pilot is over and we have usage.
|
||||
|
||||
- type: textarea
|
||||
id: others
|
||||
validations:
|
||||
required: true
|
||||
attributes:
|
||||
label: What other VAEs you would like to see if the pilot goes well?
|
||||
description: |
|
||||
Provide a list of the VAEs you would like to see in the future if the pilot goes well.
|
||||
|
||||
- type: textarea
|
||||
id: additional-info
|
||||
attributes:
|
||||
label: Notify the members of the team
|
||||
description: |
|
||||
Tag the following folks when submitting this feedback: @hlky @sayakpaul
|
||||
1
.github/workflows/benchmark.yml
vendored
1
.github/workflows/benchmark.yml
vendored
@@ -38,6 +38,7 @@ jobs:
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
python -m uv pip install -e [quality,test]
|
||||
python -m uv pip install pandas peft
|
||||
python -m uv pip uninstall transformers && python -m uv pip install transformers==4.48.0
|
||||
- name: Environment
|
||||
run: |
|
||||
python utils/print_env.py
|
||||
|
||||
15
.github/workflows/nightly_tests.yml
vendored
15
.github/workflows/nightly_tests.yml
vendored
@@ -265,7 +265,7 @@ jobs:
|
||||
|
||||
- name: Run PyTorch CUDA tests
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
|
||||
CUBLAS_WORKSPACE_CONFIG: :16:8
|
||||
run: |
|
||||
@@ -414,10 +414,16 @@ jobs:
|
||||
config:
|
||||
- backend: "bitsandbytes"
|
||||
test_location: "bnb"
|
||||
additional_deps: ["peft"]
|
||||
- backend: "gguf"
|
||||
test_location: "gguf"
|
||||
additional_deps: ["peft"]
|
||||
- backend: "torchao"
|
||||
test_location: "torchao"
|
||||
additional_deps: []
|
||||
- backend: "optimum_quanto"
|
||||
test_location: "quanto"
|
||||
additional_deps: []
|
||||
runs-on:
|
||||
group: aws-g6e-xlarge-plus
|
||||
container:
|
||||
@@ -435,6 +441,9 @@ jobs:
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
python -m uv pip install -e [quality,test]
|
||||
python -m uv pip install -U ${{ matrix.config.backend }}
|
||||
if [ "${{ join(matrix.config.additional_deps, ' ') }}" != "" ]; then
|
||||
python -m uv pip install ${{ join(matrix.config.additional_deps, ' ') }}
|
||||
fi
|
||||
python -m uv pip install pytest-reportlog
|
||||
- name: Environment
|
||||
run: |
|
||||
@@ -505,7 +514,7 @@ jobs:
|
||||
# shell: arch -arch arm64 bash {0}
|
||||
# env:
|
||||
# HF_HOME: /System/Volumes/Data/mnt/cache
|
||||
# HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
# HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
# run: |
|
||||
# ${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
|
||||
# --report-log=tests_torch_mps.log \
|
||||
@@ -561,7 +570,7 @@ jobs:
|
||||
# shell: arch -arch arm64 bash {0}
|
||||
# env:
|
||||
# HF_HOME: /System/Volumes/Data/mnt/cache
|
||||
# HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
# HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
# run: |
|
||||
# ${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
|
||||
# --report-log=tests_torch_mps.log \
|
||||
|
||||
17
.github/workflows/pr_style_bot.yml
vendored
Normal file
17
.github/workflows/pr_style_bot.yml
vendored
Normal file
@@ -0,0 +1,17 @@
|
||||
name: PR Style Bot
|
||||
|
||||
on:
|
||||
issue_comment:
|
||||
types: [created]
|
||||
|
||||
permissions:
|
||||
contents: write
|
||||
pull-requests: write
|
||||
|
||||
jobs:
|
||||
style:
|
||||
uses: huggingface/huggingface_hub/.github/workflows/style-bot-action.yml@main
|
||||
with:
|
||||
python_quality_dependencies: "[quality]"
|
||||
secrets:
|
||||
bot_token: ${{ secrets.GITHUB_TOKEN }}
|
||||
7
.github/workflows/pr_tests.yml
vendored
7
.github/workflows/pr_tests.yml
vendored
@@ -2,8 +2,7 @@ name: Fast tests for PRs
|
||||
|
||||
on:
|
||||
pull_request:
|
||||
branches:
|
||||
- main
|
||||
branches: [main]
|
||||
paths:
|
||||
- "src/diffusers/**.py"
|
||||
- "benchmarks/**.py"
|
||||
@@ -64,6 +63,7 @@ jobs:
|
||||
run: |
|
||||
python utils/check_copies.py
|
||||
python utils/check_dummies.py
|
||||
python utils/check_support_list.py
|
||||
make deps_table_check_updated
|
||||
- name: Check if failure
|
||||
if: ${{ failure() }}
|
||||
@@ -120,7 +120,8 @@ jobs:
|
||||
run: |
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
python -m uv pip install -e [quality,test]
|
||||
python -m uv pip install accelerate
|
||||
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
|
||||
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git --no-deps
|
||||
|
||||
- name: Environment
|
||||
run: |
|
||||
|
||||
296
.github/workflows/pr_tests_gpu.yml
vendored
Normal file
296
.github/workflows/pr_tests_gpu.yml
vendored
Normal file
@@ -0,0 +1,296 @@
|
||||
name: Fast GPU Tests on PR
|
||||
|
||||
on:
|
||||
pull_request:
|
||||
branches: main
|
||||
paths:
|
||||
- "src/diffusers/models/modeling_utils.py"
|
||||
- "src/diffusers/models/model_loading_utils.py"
|
||||
- "src/diffusers/pipelines/pipeline_utils.py"
|
||||
- "src/diffusers/pipeline_loading_utils.py"
|
||||
- "src/diffusers/loaders/lora_base.py"
|
||||
- "src/diffusers/loaders/lora_pipeline.py"
|
||||
- "src/diffusers/loaders/peft.py"
|
||||
- "tests/pipelines/test_pipelines_common.py"
|
||||
- "tests/models/test_modeling_common.py"
|
||||
workflow_dispatch:
|
||||
|
||||
concurrency:
|
||||
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
|
||||
cancel-in-progress: true
|
||||
|
||||
env:
|
||||
DIFFUSERS_IS_CI: yes
|
||||
OMP_NUM_THREADS: 8
|
||||
MKL_NUM_THREADS: 8
|
||||
HF_HUB_ENABLE_HF_TRANSFER: 1
|
||||
PYTEST_TIMEOUT: 600
|
||||
PIPELINE_USAGE_CUTOFF: 1000000000 # set high cutoff so that only always-test pipelines run
|
||||
|
||||
jobs:
|
||||
check_code_quality:
|
||||
runs-on: ubuntu-22.04
|
||||
steps:
|
||||
- uses: actions/checkout@v3
|
||||
- name: Set up Python
|
||||
uses: actions/setup-python@v4
|
||||
with:
|
||||
python-version: "3.8"
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
python -m pip install --upgrade pip
|
||||
pip install .[quality]
|
||||
- name: Check quality
|
||||
run: make quality
|
||||
- name: Check if failure
|
||||
if: ${{ failure() }}
|
||||
run: |
|
||||
echo "Quality check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make style && make quality'" >> $GITHUB_STEP_SUMMARY
|
||||
|
||||
check_repository_consistency:
|
||||
needs: check_code_quality
|
||||
runs-on: ubuntu-22.04
|
||||
steps:
|
||||
- uses: actions/checkout@v3
|
||||
- name: Set up Python
|
||||
uses: actions/setup-python@v4
|
||||
with:
|
||||
python-version: "3.8"
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
python -m pip install --upgrade pip
|
||||
pip install .[quality]
|
||||
- name: Check repo consistency
|
||||
run: |
|
||||
python utils/check_copies.py
|
||||
python utils/check_dummies.py
|
||||
python utils/check_support_list.py
|
||||
make deps_table_check_updated
|
||||
- name: Check if failure
|
||||
if: ${{ failure() }}
|
||||
run: |
|
||||
echo "Repo consistency check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make fix-copies'" >> $GITHUB_STEP_SUMMARY
|
||||
|
||||
setup_torch_cuda_pipeline_matrix:
|
||||
needs: [check_code_quality, check_repository_consistency]
|
||||
name: Setup Torch Pipelines CUDA Slow Tests Matrix
|
||||
runs-on:
|
||||
group: aws-general-8-plus
|
||||
container:
|
||||
image: diffusers/diffusers-pytorch-cpu
|
||||
outputs:
|
||||
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
uses: actions/checkout@v3
|
||||
with:
|
||||
fetch-depth: 2
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
python -m uv pip install -e [quality,test]
|
||||
- name: Environment
|
||||
run: |
|
||||
python utils/print_env.py
|
||||
- name: Fetch Pipeline Matrix
|
||||
id: fetch_pipeline_matrix
|
||||
run: |
|
||||
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
|
||||
echo $matrix
|
||||
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
|
||||
- name: Pipeline Tests Artifacts
|
||||
if: ${{ always() }}
|
||||
uses: actions/upload-artifact@v4
|
||||
with:
|
||||
name: test-pipelines.json
|
||||
path: reports
|
||||
|
||||
torch_pipelines_cuda_tests:
|
||||
name: Torch Pipelines CUDA Tests
|
||||
needs: setup_torch_cuda_pipeline_matrix
|
||||
strategy:
|
||||
fail-fast: false
|
||||
max-parallel: 8
|
||||
matrix:
|
||||
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
|
||||
runs-on:
|
||||
group: aws-g4dn-2xlarge
|
||||
container:
|
||||
image: diffusers/diffusers-pytorch-cuda
|
||||
options: --shm-size "16gb" --ipc host --gpus 0
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
uses: actions/checkout@v3
|
||||
with:
|
||||
fetch-depth: 2
|
||||
|
||||
- name: NVIDIA-SMI
|
||||
run: |
|
||||
nvidia-smi
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
python -m uv pip install -e [quality,test]
|
||||
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
|
||||
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
|
||||
|
||||
- name: Environment
|
||||
run: |
|
||||
python utils/print_env.py
|
||||
- name: Extract tests
|
||||
id: extract_tests
|
||||
run: |
|
||||
pattern=$(python utils/extract_tests_from_mixin.py --type pipeline)
|
||||
echo "$pattern" > /tmp/test_pattern.txt
|
||||
echo "pattern_file=/tmp/test_pattern.txt" >> $GITHUB_OUTPUT
|
||||
|
||||
- name: PyTorch CUDA checkpoint tests on Ubuntu
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
|
||||
CUBLAS_WORKSPACE_CONFIG: :16:8
|
||||
run: |
|
||||
if [ "${{ matrix.module }}" = "ip_adapters" ]; then
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "not Flax and not Onnx" \
|
||||
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
|
||||
tests/pipelines/${{ matrix.module }}
|
||||
else
|
||||
pattern=$(cat ${{ steps.extract_tests.outputs.pattern_file }})
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "not Flax and not Onnx and $pattern" \
|
||||
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
|
||||
tests/pipelines/${{ matrix.module }}
|
||||
fi
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
run: |
|
||||
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
|
||||
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
|
||||
- name: Test suite reports artifacts
|
||||
if: ${{ always() }}
|
||||
uses: actions/upload-artifact@v4
|
||||
with:
|
||||
name: pipeline_${{ matrix.module }}_test_reports
|
||||
path: reports
|
||||
|
||||
torch_cuda_tests:
|
||||
name: Torch CUDA Tests
|
||||
needs: [check_code_quality, check_repository_consistency]
|
||||
runs-on:
|
||||
group: aws-g4dn-2xlarge
|
||||
container:
|
||||
image: diffusers/diffusers-pytorch-cuda
|
||||
options: --shm-size "16gb" --ipc host --gpus 0
|
||||
defaults:
|
||||
run:
|
||||
shell: bash
|
||||
strategy:
|
||||
fail-fast: false
|
||||
max-parallel: 2
|
||||
matrix:
|
||||
module: [models, schedulers, lora, others]
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
uses: actions/checkout@v3
|
||||
with:
|
||||
fetch-depth: 2
|
||||
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
python -m uv pip install -e [quality,test]
|
||||
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
|
||||
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
|
||||
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
|
||||
|
||||
- name: Environment
|
||||
run: |
|
||||
python utils/print_env.py
|
||||
|
||||
- name: Extract tests
|
||||
id: extract_tests
|
||||
run: |
|
||||
pattern=$(python utils/extract_tests_from_mixin.py --type ${{ matrix.module }})
|
||||
echo "$pattern" > /tmp/test_pattern.txt
|
||||
echo "pattern_file=/tmp/test_pattern.txt" >> $GITHUB_OUTPUT
|
||||
|
||||
- name: Run PyTorch CUDA tests
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
|
||||
CUBLAS_WORKSPACE_CONFIG: :16:8
|
||||
run: |
|
||||
pattern=$(cat ${{ steps.extract_tests.outputs.pattern_file }})
|
||||
if [ -z "$pattern" ]; then
|
||||
python -m pytest -n 1 -sv --max-worker-restart=0 --dist=loadfile -k "not Flax and not Onnx" tests/${{ matrix.module }} \
|
||||
--make-reports=tests_torch_cuda_${{ matrix.module }}
|
||||
else
|
||||
python -m pytest -n 1 -sv --max-worker-restart=0 --dist=loadfile -k "not Flax and not Onnx and $pattern" tests/${{ matrix.module }} \
|
||||
--make-reports=tests_torch_cuda_${{ matrix.module }}
|
||||
fi
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
run: |
|
||||
cat reports/tests_torch_cuda_${{ matrix.module }}_stats.txt
|
||||
cat reports/tests_torch_cuda_${{ matrix.module }}_failures_short.txt
|
||||
|
||||
- name: Test suite reports artifacts
|
||||
if: ${{ always() }}
|
||||
uses: actions/upload-artifact@v4
|
||||
with:
|
||||
name: torch_cuda_test_reports_${{ matrix.module }}
|
||||
path: reports
|
||||
|
||||
run_examples_tests:
|
||||
name: Examples PyTorch CUDA tests on Ubuntu
|
||||
needs: [check_code_quality, check_repository_consistency]
|
||||
runs-on:
|
||||
group: aws-g4dn-2xlarge
|
||||
|
||||
container:
|
||||
image: diffusers/diffusers-pytorch-cuda
|
||||
options: --gpus 0 --shm-size "16gb" --ipc host
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
uses: actions/checkout@v3
|
||||
with:
|
||||
fetch-depth: 2
|
||||
|
||||
- name: NVIDIA-SMI
|
||||
run: |
|
||||
nvidia-smi
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
|
||||
python -m uv pip install -e [quality,test,training]
|
||||
|
||||
- name: Environment
|
||||
run: |
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
python utils/print_env.py
|
||||
|
||||
- name: Run example tests on GPU
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
run: |
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
python -m uv pip install timm
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=examples_torch_cuda examples/
|
||||
|
||||
- name: Failure short reports
|
||||
if: ${{ failure() }}
|
||||
run: |
|
||||
cat reports/examples_torch_cuda_stats.txt
|
||||
cat reports/examples_torch_cuda_failures_short.txt
|
||||
|
||||
- name: Test suite reports artifacts
|
||||
if: ${{ always() }}
|
||||
uses: actions/upload-artifact@v4
|
||||
with:
|
||||
name: examples_test_reports
|
||||
path: reports
|
||||
|
||||
12
.github/workflows/push_tests.yml
vendored
12
.github/workflows/push_tests.yml
vendored
@@ -187,7 +187,7 @@ jobs:
|
||||
|
||||
- name: Run Flax TPU tests
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 0 \
|
||||
-s -v -k "Flax" \
|
||||
@@ -235,7 +235,7 @@ jobs:
|
||||
|
||||
- name: Run ONNXRuntime CUDA tests
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "Onnx" \
|
||||
@@ -283,7 +283,7 @@ jobs:
|
||||
python utils/print_env.py
|
||||
- name: Run example tests on GPU
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
RUN_COMPILE: yes
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "compile" --make-reports=tests_torch_compile_cuda tests/
|
||||
@@ -326,7 +326,7 @@ jobs:
|
||||
python utils/print_env.py
|
||||
- name: Run example tests on GPU
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "xformers" --make-reports=tests_torch_xformers_cuda tests/
|
||||
- name: Failure short reports
|
||||
@@ -349,7 +349,6 @@ jobs:
|
||||
container:
|
||||
image: diffusers/diffusers-pytorch-cuda
|
||||
options: --gpus 0 --shm-size "16gb" --ipc host
|
||||
|
||||
steps:
|
||||
- name: Checkout diffusers
|
||||
uses: actions/checkout@v3
|
||||
@@ -359,7 +358,6 @@ jobs:
|
||||
- name: NVIDIA-SMI
|
||||
run: |
|
||||
nvidia-smi
|
||||
|
||||
- name: Install dependencies
|
||||
run: |
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
@@ -372,7 +370,7 @@ jobs:
|
||||
|
||||
- name: Run example tests on GPU
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
run: |
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
python -m uv pip install timm
|
||||
|
||||
16
.github/workflows/release_tests_fast.yml
vendored
16
.github/workflows/release_tests_fast.yml
vendored
@@ -81,7 +81,7 @@ jobs:
|
||||
python utils/print_env.py
|
||||
- name: Slow PyTorch CUDA checkpoint tests on Ubuntu
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
|
||||
CUBLAS_WORKSPACE_CONFIG: :16:8
|
||||
run: |
|
||||
@@ -135,7 +135,7 @@ jobs:
|
||||
|
||||
- name: Run PyTorch CUDA tests
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
|
||||
CUBLAS_WORKSPACE_CONFIG: :16:8
|
||||
run: |
|
||||
@@ -186,7 +186,7 @@ jobs:
|
||||
|
||||
- name: Run PyTorch CUDA tests
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
|
||||
CUBLAS_WORKSPACE_CONFIG: :16:8
|
||||
run: |
|
||||
@@ -241,7 +241,7 @@ jobs:
|
||||
|
||||
- name: Run slow Flax TPU tests
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 0 \
|
||||
-s -v -k "Flax" \
|
||||
@@ -289,7 +289,7 @@ jobs:
|
||||
|
||||
- name: Run slow ONNXRuntime CUDA tests
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
|
||||
-s -v -k "Onnx" \
|
||||
@@ -337,7 +337,7 @@ jobs:
|
||||
python utils/print_env.py
|
||||
- name: Run example tests on GPU
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
RUN_COMPILE: yes
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "compile" --make-reports=tests_torch_compile_cuda tests/
|
||||
@@ -380,7 +380,7 @@ jobs:
|
||||
python utils/print_env.py
|
||||
- name: Run example tests on GPU
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
run: |
|
||||
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "xformers" --make-reports=tests_torch_xformers_cuda tests/
|
||||
- name: Failure short reports
|
||||
@@ -426,7 +426,7 @@ jobs:
|
||||
|
||||
- name: Run example tests on GPU
|
||||
env:
|
||||
HF_TOKEN: ${{ secrets.HF_TOKEN }}
|
||||
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
|
||||
run: |
|
||||
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
|
||||
python -m uv pip install timm
|
||||
|
||||
14
.github/workflows/run_tests_from_a_pr.yml
vendored
14
.github/workflows/run_tests_from_a_pr.yml
vendored
@@ -7,8 +7,8 @@ on:
|
||||
default: 'diffusers/diffusers-pytorch-cuda'
|
||||
description: 'Name of the Docker image'
|
||||
required: true
|
||||
branch:
|
||||
description: 'PR Branch to test on'
|
||||
pr_number:
|
||||
description: 'PR number to test on'
|
||||
required: true
|
||||
test:
|
||||
description: 'Tests to run (e.g.: `tests/models`).'
|
||||
@@ -43,8 +43,8 @@ jobs:
|
||||
exit 1
|
||||
fi
|
||||
|
||||
if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines) ]]; then
|
||||
echo "Error: The input string must contain either 'models' or 'pipelines' after 'tests/'."
|
||||
if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines|lora) ]]; then
|
||||
echo "Error: The input string must contain either 'models', 'pipelines', or 'lora' after 'tests/'."
|
||||
exit 1
|
||||
fi
|
||||
|
||||
@@ -53,13 +53,13 @@ jobs:
|
||||
exit 1
|
||||
fi
|
||||
echo "$PY_TEST"
|
||||
|
||||
shell: bash -e {0}
|
||||
|
||||
- name: Checkout PR branch
|
||||
uses: actions/checkout@v4
|
||||
with:
|
||||
ref: ${{ github.event.inputs.branch }}
|
||||
repository: ${{ github.event.pull_request.head.repo.full_name }}
|
||||
|
||||
ref: refs/pull/${{ inputs.pr_number }}/head
|
||||
|
||||
- name: Install pytest
|
||||
run: |
|
||||
|
||||
3
.github/workflows/trufflehog.yml
vendored
3
.github/workflows/trufflehog.yml
vendored
@@ -13,3 +13,6 @@ jobs:
|
||||
fetch-depth: 0
|
||||
- name: Secret Scanning
|
||||
uses: trufflesecurity/trufflehog@main
|
||||
with:
|
||||
extra_args: --results=verified,unknown
|
||||
|
||||
|
||||
@@ -28,9 +28,9 @@ ENV PATH="/opt/venv/bin:$PATH"
|
||||
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
|
||||
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
|
||||
python3 -m uv pip install --no-cache-dir \
|
||||
torch==2.1.2 \
|
||||
torchvision==0.16.2 \
|
||||
torchaudio==2.1.2 \
|
||||
torch \
|
||||
torchvision \
|
||||
torchaudio\
|
||||
onnxruntime \
|
||||
--extra-index-url https://download.pytorch.org/whl/cpu && \
|
||||
python3 -m uv pip install --no-cache-dir \
|
||||
|
||||
@@ -76,6 +76,16 @@
|
||||
- local: advanced_inference/outpaint
|
||||
title: Outpainting
|
||||
title: Advanced inference
|
||||
- sections:
|
||||
- local: hybrid_inference/overview
|
||||
title: Overview
|
||||
- local: hybrid_inference/vae_decode
|
||||
title: VAE Decode
|
||||
- local: hybrid_inference/vae_encode
|
||||
title: VAE Encode
|
||||
- local: hybrid_inference/api_reference
|
||||
title: API Reference
|
||||
title: Hybrid Inference
|
||||
- sections:
|
||||
- local: using-diffusers/cogvideox
|
||||
title: CogVideoX
|
||||
@@ -89,6 +99,8 @@
|
||||
title: Kandinsky
|
||||
- local: using-diffusers/ip_adapter
|
||||
title: IP-Adapter
|
||||
- local: using-diffusers/omnigen
|
||||
title: OmniGen
|
||||
- local: using-diffusers/pag
|
||||
title: PAG
|
||||
- local: using-diffusers/controlnet
|
||||
@@ -163,6 +175,8 @@
|
||||
title: gguf
|
||||
- local: quantization/torchao
|
||||
title: torchao
|
||||
- local: quantization/quanto
|
||||
title: quanto
|
||||
title: Quantization Methods
|
||||
- sections:
|
||||
- local: optimization/fp16
|
||||
@@ -251,19 +265,23 @@
|
||||
sections:
|
||||
- local: api/models/overview
|
||||
title: Overview
|
||||
- local: api/models/auto_model
|
||||
title: AutoModel
|
||||
- sections:
|
||||
- local: api/models/controlnet
|
||||
title: ControlNetModel
|
||||
- local: api/models/controlnet_union
|
||||
title: ControlNetUnionModel
|
||||
- local: api/models/controlnet_flux
|
||||
title: FluxControlNetModel
|
||||
- local: api/models/controlnet_hunyuandit
|
||||
title: HunyuanDiT2DControlNetModel
|
||||
- local: api/models/controlnet_sana
|
||||
title: SanaControlNetModel
|
||||
- local: api/models/controlnet_sd3
|
||||
title: SD3ControlNetModel
|
||||
- local: api/models/controlnet_sparsectrl
|
||||
title: SparseControlNetModel
|
||||
- local: api/models/controlnet_union
|
||||
title: ControlNetUnionModel
|
||||
title: ControlNets
|
||||
- sections:
|
||||
- local: api/models/allegro_transformer3d
|
||||
@@ -272,50 +290,62 @@
|
||||
title: AuraFlowTransformer2DModel
|
||||
- local: api/models/cogvideox_transformer3d
|
||||
title: CogVideoXTransformer3DModel
|
||||
- local: api/models/consisid_transformer3d
|
||||
title: ConsisIDTransformer3DModel
|
||||
- local: api/models/cogview3plus_transformer2d
|
||||
title: CogView3PlusTransformer2DModel
|
||||
- local: api/models/cogview4_transformer2d
|
||||
title: CogView4Transformer2DModel
|
||||
- local: api/models/consisid_transformer3d
|
||||
title: ConsisIDTransformer3DModel
|
||||
- local: api/models/dit_transformer2d
|
||||
title: DiTTransformer2DModel
|
||||
- local: api/models/easyanimate_transformer3d
|
||||
title: EasyAnimateTransformer3DModel
|
||||
- local: api/models/flux_transformer
|
||||
title: FluxTransformer2DModel
|
||||
- local: api/models/hidream_image_transformer
|
||||
title: HiDreamImageTransformer2DModel
|
||||
- local: api/models/hunyuan_transformer2d
|
||||
title: HunyuanDiT2DModel
|
||||
- local: api/models/hunyuan_video_transformer_3d
|
||||
title: HunyuanVideoTransformer3DModel
|
||||
- local: api/models/latte_transformer3d
|
||||
title: LatteTransformer3DModel
|
||||
- local: api/models/lumina_nextdit2d
|
||||
title: LuminaNextDiT2DModel
|
||||
- local: api/models/ltx_video_transformer3d
|
||||
title: LTXVideoTransformer3DModel
|
||||
- local: api/models/lumina2_transformer2d
|
||||
title: Lumina2Transformer2DModel
|
||||
- local: api/models/lumina_nextdit2d
|
||||
title: LuminaNextDiT2DModel
|
||||
- local: api/models/mochi_transformer3d
|
||||
title: MochiTransformer3DModel
|
||||
- local: api/models/omnigen_transformer
|
||||
title: OmniGenTransformer2DModel
|
||||
- local: api/models/pixart_transformer2d
|
||||
title: PixArtTransformer2DModel
|
||||
- local: api/models/prior_transformer
|
||||
title: PriorTransformer
|
||||
- local: api/models/sd3_transformer2d
|
||||
title: SD3Transformer2DModel
|
||||
- local: api/models/sana_transformer2d
|
||||
title: SanaTransformer2DModel
|
||||
- local: api/models/sd3_transformer2d
|
||||
title: SD3Transformer2DModel
|
||||
- local: api/models/stable_audio_transformer
|
||||
title: StableAudioDiTModel
|
||||
- local: api/models/transformer2d
|
||||
title: Transformer2DModel
|
||||
- local: api/models/transformer_temporal
|
||||
title: TransformerTemporalModel
|
||||
- local: api/models/wan_transformer_3d
|
||||
title: WanTransformer3DModel
|
||||
title: Transformers
|
||||
- sections:
|
||||
- local: api/models/stable_cascade_unet
|
||||
title: StableCascadeUNet
|
||||
- local: api/models/unet
|
||||
title: UNet1DModel
|
||||
- local: api/models/unet2d
|
||||
title: UNet2DModel
|
||||
- local: api/models/unet2d-cond
|
||||
title: UNet2DConditionModel
|
||||
- local: api/models/unet2d
|
||||
title: UNet2DModel
|
||||
- local: api/models/unet3d-cond
|
||||
title: UNet3DConditionModel
|
||||
- local: api/models/unet-motion
|
||||
@@ -324,6 +354,10 @@
|
||||
title: UViT2DModel
|
||||
title: UNets
|
||||
- sections:
|
||||
- local: api/models/asymmetricautoencoderkl
|
||||
title: AsymmetricAutoencoderKL
|
||||
- local: api/models/autoencoder_dc
|
||||
title: AutoencoderDC
|
||||
- local: api/models/autoencoderkl
|
||||
title: AutoencoderKL
|
||||
- local: api/models/autoencoderkl_allegro
|
||||
@@ -334,12 +368,12 @@
|
||||
title: AutoencoderKLHunyuanVideo
|
||||
- local: api/models/autoencoderkl_ltx_video
|
||||
title: AutoencoderKLLTXVideo
|
||||
- local: api/models/autoencoderkl_magvit
|
||||
title: AutoencoderKLMagvit
|
||||
- local: api/models/autoencoderkl_mochi
|
||||
title: AutoencoderKLMochi
|
||||
- local: api/models/asymmetricautoencoderkl
|
||||
title: AsymmetricAutoencoderKL
|
||||
- local: api/models/autoencoder_dc
|
||||
title: AutoencoderDC
|
||||
- local: api/models/autoencoder_kl_wan
|
||||
title: AutoencoderKLWan
|
||||
- local: api/models/consistency_decoder_vae
|
||||
title: ConsistencyDecoderVAE
|
||||
- local: api/models/autoencoder_oobleck
|
||||
@@ -376,6 +410,8 @@
|
||||
title: CogVideoX
|
||||
- local: api/pipelines/cogview3
|
||||
title: CogView3
|
||||
- local: api/pipelines/cogview4
|
||||
title: CogView4
|
||||
- local: api/pipelines/consisid
|
||||
title: ConsisID
|
||||
- local: api/pipelines/consistency_models
|
||||
@@ -390,6 +426,8 @@
|
||||
title: ControlNet with Stable Diffusion 3
|
||||
- local: api/pipelines/controlnet_sdxl
|
||||
title: ControlNet with Stable Diffusion XL
|
||||
- local: api/pipelines/controlnet_sana
|
||||
title: ControlNet-Sana
|
||||
- local: api/pipelines/controlnetxs
|
||||
title: ControlNet-XS
|
||||
- local: api/pipelines/controlnetxs_sdxl
|
||||
@@ -408,10 +446,14 @@
|
||||
title: DiffEdit
|
||||
- local: api/pipelines/dit
|
||||
title: DiT
|
||||
- local: api/pipelines/easyanimate
|
||||
title: EasyAnimate
|
||||
- local: api/pipelines/flux
|
||||
title: Flux
|
||||
- local: api/pipelines/control_flux_inpaint
|
||||
title: FluxControlInpaint
|
||||
- local: api/pipelines/hidream
|
||||
title: HiDream-I1
|
||||
- local: api/pipelines/hunyuandit
|
||||
title: Hunyuan-DiT
|
||||
- local: api/pipelines/hunyuan_video
|
||||
@@ -438,6 +480,8 @@
|
||||
title: LEDITS++
|
||||
- local: api/pipelines/ltx_video
|
||||
title: LTXVideo
|
||||
- local: api/pipelines/lumina2
|
||||
title: Lumina 2.0
|
||||
- local: api/pipelines/lumina
|
||||
title: Lumina-T2X
|
||||
- local: api/pipelines/marigold
|
||||
@@ -448,6 +492,8 @@
|
||||
title: MultiDiffusion
|
||||
- local: api/pipelines/musicldm
|
||||
title: MusicLDM
|
||||
- local: api/pipelines/omnigen
|
||||
title: OmniGen
|
||||
- local: api/pipelines/pag
|
||||
title: PAG
|
||||
- local: api/pipelines/paint_by_example
|
||||
@@ -460,6 +506,8 @@
|
||||
title: PixArt-Σ
|
||||
- local: api/pipelines/sana
|
||||
title: Sana
|
||||
- local: api/pipelines/sana_sprint
|
||||
title: Sana Sprint
|
||||
- local: api/pipelines/self_attention_guidance
|
||||
title: Self-Attention Guidance
|
||||
- local: api/pipelines/semantic_stable_diffusion
|
||||
@@ -473,40 +521,40 @@
|
||||
- sections:
|
||||
- local: api/pipelines/stable_diffusion/overview
|
||||
title: Overview
|
||||
- local: api/pipelines/stable_diffusion/text2img
|
||||
title: Text-to-image
|
||||
- local: api/pipelines/stable_diffusion/depth2img
|
||||
title: Depth-to-image
|
||||
- local: api/pipelines/stable_diffusion/gligen
|
||||
title: GLIGEN (Grounded Language-to-Image Generation)
|
||||
- local: api/pipelines/stable_diffusion/image_variation
|
||||
title: Image variation
|
||||
- local: api/pipelines/stable_diffusion/img2img
|
||||
title: Image-to-image
|
||||
- local: api/pipelines/stable_diffusion/svd
|
||||
title: Image-to-video
|
||||
- local: api/pipelines/stable_diffusion/inpaint
|
||||
title: Inpainting
|
||||
- local: api/pipelines/stable_diffusion/depth2img
|
||||
title: Depth-to-image
|
||||
- local: api/pipelines/stable_diffusion/image_variation
|
||||
title: Image variation
|
||||
- local: api/pipelines/stable_diffusion/k_diffusion
|
||||
title: K-Diffusion
|
||||
- local: api/pipelines/stable_diffusion/latent_upscale
|
||||
title: Latent upscaler
|
||||
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
|
||||
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D Upscaler
|
||||
- local: api/pipelines/stable_diffusion/stable_diffusion_safe
|
||||
title: Safe Stable Diffusion
|
||||
- local: api/pipelines/stable_diffusion/sdxl_turbo
|
||||
title: SDXL Turbo
|
||||
- local: api/pipelines/stable_diffusion/stable_diffusion_2
|
||||
title: Stable Diffusion 2
|
||||
- local: api/pipelines/stable_diffusion/stable_diffusion_3
|
||||
title: Stable Diffusion 3
|
||||
- local: api/pipelines/stable_diffusion/stable_diffusion_xl
|
||||
title: Stable Diffusion XL
|
||||
- local: api/pipelines/stable_diffusion/sdxl_turbo
|
||||
title: SDXL Turbo
|
||||
- local: api/pipelines/stable_diffusion/latent_upscale
|
||||
title: Latent upscaler
|
||||
- local: api/pipelines/stable_diffusion/upscale
|
||||
title: Super-resolution
|
||||
- local: api/pipelines/stable_diffusion/k_diffusion
|
||||
title: K-Diffusion
|
||||
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
|
||||
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D Upscaler
|
||||
- local: api/pipelines/stable_diffusion/adapter
|
||||
title: T2I-Adapter
|
||||
- local: api/pipelines/stable_diffusion/gligen
|
||||
title: GLIGEN (Grounded Language-to-Image Generation)
|
||||
- local: api/pipelines/stable_diffusion/text2img
|
||||
title: Text-to-image
|
||||
title: Stable Diffusion
|
||||
- local: api/pipelines/stable_unclip
|
||||
title: Stable unCLIP
|
||||
@@ -520,6 +568,8 @@
|
||||
title: UniDiffuser
|
||||
- local: api/pipelines/value_guided_sampling
|
||||
title: Value-guided sampling
|
||||
- local: api/pipelines/wan
|
||||
title: Wan
|
||||
- local: api/pipelines/wuerstchen
|
||||
title: Wuerstchen
|
||||
title: Pipelines
|
||||
@@ -529,6 +579,10 @@
|
||||
title: Overview
|
||||
- local: api/schedulers/cm_stochastic_iterative
|
||||
title: CMStochasticIterativeScheduler
|
||||
- local: api/schedulers/ddim_cogvideox
|
||||
title: CogVideoXDDIMScheduler
|
||||
- local: api/schedulers/multistep_dpm_solver_cogvideox
|
||||
title: CogVideoXDPMScheduler
|
||||
- local: api/schedulers/consistency_decoder
|
||||
title: ConsistencyDecoderScheduler
|
||||
- local: api/schedulers/cosine_dpm
|
||||
@@ -598,6 +652,8 @@
|
||||
title: Attention Processor
|
||||
- local: api/activations
|
||||
title: Custom activation functions
|
||||
- local: api/cache
|
||||
title: Caching methods
|
||||
- local: api/normalization
|
||||
title: Custom normalization layers
|
||||
- local: api/utilities
|
||||
|
||||
@@ -25,3 +25,16 @@ Customized activation functions for supporting various models in 🤗 Diffusers.
|
||||
## ApproximateGELU
|
||||
|
||||
[[autodoc]] models.activations.ApproximateGELU
|
||||
|
||||
|
||||
## SwiGLU
|
||||
|
||||
[[autodoc]] models.activations.SwiGLU
|
||||
|
||||
## FP32SiLU
|
||||
|
||||
[[autodoc]] models.activations.FP32SiLU
|
||||
|
||||
## LinearActivation
|
||||
|
||||
[[autodoc]] models.activations.LinearActivation
|
||||
|
||||
@@ -147,3 +147,20 @@ An attention processor is a class for applying different types of attention mech
|
||||
## XLAFlashAttnProcessor2_0
|
||||
|
||||
[[autodoc]] models.attention_processor.XLAFlashAttnProcessor2_0
|
||||
|
||||
## XFormersJointAttnProcessor
|
||||
|
||||
[[autodoc]] models.attention_processor.XFormersJointAttnProcessor
|
||||
|
||||
## IPAdapterXFormersAttnProcessor
|
||||
|
||||
[[autodoc]] models.attention_processor.IPAdapterXFormersAttnProcessor
|
||||
|
||||
## FluxIPAdapterJointAttnProcessor2_0
|
||||
|
||||
[[autodoc]] models.attention_processor.FluxIPAdapterJointAttnProcessor2_0
|
||||
|
||||
|
||||
## XLAFluxFlashAttnProcessor2_0
|
||||
|
||||
[[autodoc]] models.attention_processor.XLAFluxFlashAttnProcessor2_0
|
||||
82
docs/source/en/api/cache.md
Normal file
82
docs/source/en/api/cache.md
Normal file
@@ -0,0 +1,82 @@
|
||||
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License. -->
|
||||
|
||||
# Caching methods
|
||||
|
||||
## Pyramid Attention Broadcast
|
||||
|
||||
[Pyramid Attention Broadcast](https://huggingface.co/papers/2408.12588) from Xuanlei Zhao, Xiaolong Jin, Kai Wang, Yang You.
|
||||
|
||||
Pyramid Attention Broadcast (PAB) is a method that speeds up inference in diffusion models by systematically skipping attention computations between successive inference steps and reusing cached attention states. The attention states are not very different between successive inference steps. The most prominent difference is in the spatial attention blocks, not as much in the temporal attention blocks, and finally the least in the cross attention blocks. Therefore, many cross attention computation blocks can be skipped, followed by the temporal and spatial attention blocks. By combining other techniques like sequence parallelism and classifier-free guidance parallelism, PAB achieves near real-time video generation.
|
||||
|
||||
Enable PAB with [`~PyramidAttentionBroadcastConfig`] on any pipeline. For some benchmarks, refer to [this](https://github.com/huggingface/diffusers/pull/9562) pull request.
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import CogVideoXPipeline, PyramidAttentionBroadcastConfig
|
||||
|
||||
pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
|
||||
pipe.to("cuda")
|
||||
|
||||
# Increasing the value of `spatial_attention_timestep_skip_range[0]` or decreasing the value of
|
||||
# `spatial_attention_timestep_skip_range[1]` will decrease the interval in which pyramid attention
|
||||
# broadcast is active, leader to slower inference speeds. However, large intervals can lead to
|
||||
# poorer quality of generated videos.
|
||||
config = PyramidAttentionBroadcastConfig(
|
||||
spatial_attention_block_skip_range=2,
|
||||
spatial_attention_timestep_skip_range=(100, 800),
|
||||
current_timestep_callback=lambda: pipe.current_timestep,
|
||||
)
|
||||
pipe.transformer.enable_cache(config)
|
||||
```
|
||||
|
||||
## Faster Cache
|
||||
|
||||
[FasterCache](https://huggingface.co/papers/2410.19355) from Zhengyao Lv, Chenyang Si, Junhao Song, Zhenyu Yang, Yu Qiao, Ziwei Liu, Kwan-Yee K. Wong.
|
||||
|
||||
FasterCache is a method that speeds up inference in diffusion transformers by:
|
||||
- Reusing attention states between successive inference steps, due to high similarity between them
|
||||
- Skipping unconditional branch prediction used in classifier-free guidance by revealing redundancies between unconditional and conditional branch outputs for the same timestep, and therefore approximating the unconditional branch output using the conditional branch output
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import CogVideoXPipeline, FasterCacheConfig
|
||||
|
||||
pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
|
||||
pipe.to("cuda")
|
||||
|
||||
config = FasterCacheConfig(
|
||||
spatial_attention_block_skip_range=2,
|
||||
spatial_attention_timestep_skip_range=(-1, 681),
|
||||
current_timestep_callback=lambda: pipe.current_timestep,
|
||||
attention_weight_callback=lambda _: 0.3,
|
||||
unconditional_batch_skip_range=5,
|
||||
unconditional_batch_timestep_skip_range=(-1, 781),
|
||||
tensor_format="BFCHW",
|
||||
)
|
||||
pipe.transformer.enable_cache(config)
|
||||
```
|
||||
|
||||
### CacheMixin
|
||||
|
||||
[[autodoc]] CacheMixin
|
||||
|
||||
### PyramidAttentionBroadcastConfig
|
||||
|
||||
[[autodoc]] PyramidAttentionBroadcastConfig
|
||||
|
||||
[[autodoc]] apply_pyramid_attention_broadcast
|
||||
|
||||
### FasterCacheConfig
|
||||
|
||||
[[autodoc]] FasterCacheConfig
|
||||
|
||||
[[autodoc]] apply_faster_cache
|
||||
@@ -20,6 +20,13 @@ LoRA is a fast and lightweight training method that inserts and trains a signifi
|
||||
- [`FluxLoraLoaderMixin`] provides similar functions for [Flux](https://huggingface.co/docs/diffusers/main/en/api/pipelines/flux).
|
||||
- [`CogVideoXLoraLoaderMixin`] provides similar functions for [CogVideoX](https://huggingface.co/docs/diffusers/main/en/api/pipelines/cogvideox).
|
||||
- [`Mochi1LoraLoaderMixin`] provides similar functions for [Mochi](https://huggingface.co/docs/diffusers/main/en/api/pipelines/mochi).
|
||||
- [`AuraFlowLoraLoaderMixin`] provides similar functions for [AuraFlow](https://huggingface.co/fal/AuraFlow).
|
||||
- [`LTXVideoLoraLoaderMixin`] provides similar functions for [LTX-Video](https://huggingface.co/docs/diffusers/main/en/api/pipelines/ltx_video).
|
||||
- [`SanaLoraLoaderMixin`] provides similar functions for [Sana](https://huggingface.co/docs/diffusers/main/en/api/pipelines/sana).
|
||||
- [`HunyuanVideoLoraLoaderMixin`] provides similar functions for [HunyuanVideo](https://huggingface.co/docs/diffusers/main/en/api/pipelines/hunyuan_video).
|
||||
- [`Lumina2LoraLoaderMixin`] provides similar functions for [Lumina2](https://huggingface.co/docs/diffusers/main/en/api/pipelines/lumina2).
|
||||
- [`WanLoraLoaderMixin`] provides similar functions for [Wan](https://huggingface.co/docs/diffusers/main/en/api/pipelines/wan).
|
||||
- [`CogView4LoraLoaderMixin`] provides similar functions for [CogView4](https://huggingface.co/docs/diffusers/main/en/api/pipelines/cogview4).
|
||||
- [`AmusedLoraLoaderMixin`] is for the [`AmusedPipeline`].
|
||||
- [`LoraBaseMixin`] provides a base class with several utility methods to fuse, unfuse, unload, LoRAs and more.
|
||||
|
||||
@@ -52,6 +59,33 @@ To learn more about how to load LoRA weights, see the [LoRA](../../using-diffuse
|
||||
## Mochi1LoraLoaderMixin
|
||||
|
||||
[[autodoc]] loaders.lora_pipeline.Mochi1LoraLoaderMixin
|
||||
## AuraFlowLoraLoaderMixin
|
||||
|
||||
[[autodoc]] loaders.lora_pipeline.AuraFlowLoraLoaderMixin
|
||||
|
||||
## LTXVideoLoraLoaderMixin
|
||||
|
||||
[[autodoc]] loaders.lora_pipeline.LTXVideoLoraLoaderMixin
|
||||
|
||||
## SanaLoraLoaderMixin
|
||||
|
||||
[[autodoc]] loaders.lora_pipeline.SanaLoraLoaderMixin
|
||||
|
||||
## HunyuanVideoLoraLoaderMixin
|
||||
|
||||
[[autodoc]] loaders.lora_pipeline.HunyuanVideoLoraLoaderMixin
|
||||
|
||||
## Lumina2LoraLoaderMixin
|
||||
|
||||
[[autodoc]] loaders.lora_pipeline.Lumina2LoraLoaderMixin
|
||||
|
||||
## CogView4LoraLoaderMixin
|
||||
|
||||
[[autodoc]] loaders.lora_pipeline.CogView4LoraLoaderMixin
|
||||
|
||||
## WanLoraLoaderMixin
|
||||
|
||||
[[autodoc]] loaders.lora_pipeline.WanLoraLoaderMixin
|
||||
|
||||
## AmusedLoraLoaderMixin
|
||||
|
||||
|
||||
29
docs/source/en/api/models/auto_model.md
Normal file
29
docs/source/en/api/models/auto_model.md
Normal file
@@ -0,0 +1,29 @@
|
||||
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# AutoModel
|
||||
|
||||
The `AutoModel` is designed to make it easy to load a checkpoint without needing to know the specific model class. `AutoModel` automatically retrieves the correct model class from the checkpoint `config.json` file.
|
||||
|
||||
```python
|
||||
from diffusers import AutoModel, AutoPipelineForText2Image
|
||||
|
||||
unet = AutoModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", subfolder="unet")
|
||||
pipe = AutoPipelineForText2Image.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", unet=unet)
|
||||
```
|
||||
|
||||
|
||||
## AutoModel
|
||||
|
||||
[[autodoc]] AutoModel
|
||||
- all
|
||||
- from_pretrained
|
||||
32
docs/source/en/api/models/autoencoder_kl_wan.md
Normal file
32
docs/source/en/api/models/autoencoder_kl_wan.md
Normal file
@@ -0,0 +1,32 @@
|
||||
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License. -->
|
||||
|
||||
# AutoencoderKLWan
|
||||
|
||||
The 3D variational autoencoder (VAE) model with KL loss used in [Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
|
||||
|
||||
The model can be loaded with the following code snippet.
|
||||
|
||||
```python
|
||||
from diffusers import AutoencoderKLWan
|
||||
|
||||
vae = AutoencoderKLWan.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", subfolder="vae", torch_dtype=torch.float32)
|
||||
```
|
||||
|
||||
## AutoencoderKLWan
|
||||
|
||||
[[autodoc]] AutoencoderKLWan
|
||||
- decode
|
||||
- all
|
||||
|
||||
## DecoderOutput
|
||||
|
||||
[[autodoc]] models.autoencoders.vae.DecoderOutput
|
||||
@@ -18,7 +18,7 @@ The model can be loaded with the following code snippet.
|
||||
```python
|
||||
from diffusers import AutoencoderKLAllegro
|
||||
|
||||
vae = AutoencoderKLCogVideoX.from_pretrained("rhymes-ai/Allegro", subfolder="vae", torch_dtype=torch.float32).to("cuda")
|
||||
vae = AutoencoderKLAllegro.from_pretrained("rhymes-ai/Allegro", subfolder="vae", torch_dtype=torch.float32).to("cuda")
|
||||
```
|
||||
|
||||
## AutoencoderKLAllegro
|
||||
|
||||
37
docs/source/en/api/models/autoencoderkl_magvit.md
Normal file
37
docs/source/en/api/models/autoencoderkl_magvit.md
Normal file
@@ -0,0 +1,37 @@
|
||||
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License. -->
|
||||
|
||||
# AutoencoderKLMagvit
|
||||
|
||||
The 3D variational autoencoder (VAE) model with KL loss used in [EasyAnimate](https://github.com/aigc-apps/EasyAnimate) was introduced by Alibaba PAI.
|
||||
|
||||
The model can be loaded with the following code snippet.
|
||||
|
||||
```python
|
||||
from diffusers import AutoencoderKLMagvit
|
||||
|
||||
vae = AutoencoderKLMagvit.from_pretrained("alibaba-pai/EasyAnimateV5.1-12b-zh", subfolder="vae", torch_dtype=torch.float16).to("cuda")
|
||||
```
|
||||
|
||||
## AutoencoderKLMagvit
|
||||
|
||||
[[autodoc]] AutoencoderKLMagvit
|
||||
- decode
|
||||
- encode
|
||||
- all
|
||||
|
||||
## AutoencoderKLOutput
|
||||
|
||||
[[autodoc]] models.autoencoders.autoencoder_kl.AutoencoderKLOutput
|
||||
|
||||
## DecoderOutput
|
||||
|
||||
[[autodoc]] models.autoencoders.vae.DecoderOutput
|
||||
30
docs/source/en/api/models/cogview4_transformer2d.md
Normal file
30
docs/source/en/api/models/cogview4_transformer2d.md
Normal file
@@ -0,0 +1,30 @@
|
||||
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License. -->
|
||||
|
||||
# CogView4Transformer2DModel
|
||||
|
||||
A Diffusion Transformer model for 2D data from [CogView4]()
|
||||
|
||||
The model can be loaded with the following code snippet.
|
||||
|
||||
```python
|
||||
from diffusers import CogView4Transformer2DModel
|
||||
|
||||
transformer = CogView4Transformer2DModel.from_pretrained("THUDM/CogView4-6B", subfolder="transformer", torch_dtype=torch.bfloat16).to("cuda")
|
||||
```
|
||||
|
||||
## CogView4Transformer2DModel
|
||||
|
||||
[[autodoc]] CogView4Transformer2DModel
|
||||
|
||||
## Transformer2DModelOutput
|
||||
|
||||
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput
|
||||
29
docs/source/en/api/models/controlnet_sana.md
Normal file
29
docs/source/en/api/models/controlnet_sana.md
Normal file
@@ -0,0 +1,29 @@
|
||||
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# SanaControlNetModel
|
||||
|
||||
The ControlNet model was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, Maneesh Agrawala. It provides a greater degree of control over text-to-image generation by conditioning the model on additional inputs such as edge maps, depth maps, segmentation maps, and keypoints for pose detection.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
|
||||
|
||||
This model was contributed by [ishan24](https://huggingface.co/ishan24). ❤️
|
||||
The original codebase can be found at [NVlabs/Sana](https://github.com/NVlabs/Sana), and you can find official ControlNet checkpoints on [Efficient-Large-Model's](https://huggingface.co/Efficient-Large-Model) Hub profile.
|
||||
|
||||
## SanaControlNetModel
|
||||
[[autodoc]] SanaControlNetModel
|
||||
|
||||
## SanaControlNetOutput
|
||||
[[autodoc]] models.controlnets.controlnet_sana.SanaControlNetOutput
|
||||
|
||||
30
docs/source/en/api/models/easyanimate_transformer3d.md
Normal file
30
docs/source/en/api/models/easyanimate_transformer3d.md
Normal file
@@ -0,0 +1,30 @@
|
||||
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License. -->
|
||||
|
||||
# EasyAnimateTransformer3DModel
|
||||
|
||||
A Diffusion Transformer model for 3D data from [EasyAnimate](https://github.com/aigc-apps/EasyAnimate) was introduced by Alibaba PAI.
|
||||
|
||||
The model can be loaded with the following code snippet.
|
||||
|
||||
```python
|
||||
from diffusers import EasyAnimateTransformer3DModel
|
||||
|
||||
transformer = EasyAnimateTransformer3DModel.from_pretrained("alibaba-pai/EasyAnimateV5.1-12b-zh", subfolder="transformer", torch_dtype=torch.float16).to("cuda")
|
||||
```
|
||||
|
||||
## EasyAnimateTransformer3DModel
|
||||
|
||||
[[autodoc]] EasyAnimateTransformer3DModel
|
||||
|
||||
## Transformer2DModelOutput
|
||||
|
||||
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput
|
||||
30
docs/source/en/api/models/hidream_image_transformer.md
Normal file
30
docs/source/en/api/models/hidream_image_transformer.md
Normal file
@@ -0,0 +1,30 @@
|
||||
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License. -->
|
||||
|
||||
# HiDreamImageTransformer2DModel
|
||||
|
||||
A Transformer model for image-like data from [HiDream-I1](https://huggingface.co/HiDream-ai).
|
||||
|
||||
The model can be loaded with the following code snippet.
|
||||
|
||||
```python
|
||||
from diffusers import HiDreamImageTransformer2DModel
|
||||
|
||||
transformer = HiDreamImageTransformer2DModel.from_pretrained("HiDream-ai/HiDream-I1-Full", subfolder="transformer", torch_dtype=torch.bfloat16)
|
||||
```
|
||||
|
||||
## HiDreamImageTransformer2DModel
|
||||
|
||||
[[autodoc]] HiDreamImageTransformer2DModel
|
||||
|
||||
## Transformer2DModelOutput
|
||||
|
||||
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput
|
||||
30
docs/source/en/api/models/lumina2_transformer2d.md
Normal file
30
docs/source/en/api/models/lumina2_transformer2d.md
Normal file
@@ -0,0 +1,30 @@
|
||||
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License. -->
|
||||
|
||||
# Lumina2Transformer2DModel
|
||||
|
||||
A Diffusion Transformer model for 3D video-like data was introduced in [Lumina Image 2.0](https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0) by Alpha-VLLM.
|
||||
|
||||
The model can be loaded with the following code snippet.
|
||||
|
||||
```python
|
||||
from diffusers import Lumina2Transformer2DModel
|
||||
|
||||
transformer = Lumina2Transformer2DModel.from_pretrained("Alpha-VLLM/Lumina-Image-2.0", subfolder="transformer", torch_dtype=torch.bfloat16)
|
||||
```
|
||||
|
||||
## Lumina2Transformer2DModel
|
||||
|
||||
[[autodoc]] Lumina2Transformer2DModel
|
||||
|
||||
## Transformer2DModelOutput
|
||||
|
||||
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput
|
||||
30
docs/source/en/api/models/omnigen_transformer.md
Normal file
30
docs/source/en/api/models/omnigen_transformer.md
Normal file
@@ -0,0 +1,30 @@
|
||||
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# OmniGenTransformer2DModel
|
||||
|
||||
A Transformer model that accepts multimodal instructions to generate images for [OmniGen](https://github.com/VectorSpaceLab/OmniGen/).
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
*The emergence of Large Language Models (LLMs) has unified language generation tasks and revolutionized human-machine interaction. However, in the realm of image generation, a unified model capable of handling various tasks within a single framework remains largely unexplored. In this work, we introduce OmniGen, a new diffusion model for unified image generation. OmniGen is characterized by the following features: 1) Unification: OmniGen not only demonstrates text-to-image generation capabilities but also inherently supports various downstream tasks, such as image editing, subject-driven generation, and visual conditional generation. 2) Simplicity: The architecture of OmniGen is highly simplified, eliminating the need for additional plugins. Moreover, compared to existing diffusion models, it is more user-friendly and can complete complex tasks end-to-end through instructions without the need for extra intermediate steps, greatly simplifying the image generation workflow. 3) Knowledge Transfer: Benefit from learning in a unified format, OmniGen effectively transfers knowledge across different tasks, manages unseen tasks and domains, and exhibits novel capabilities. We also explore the model’s reasoning capabilities and potential applications of the chain-of-thought mechanism. This work represents the first attempt at a general-purpose image generation model, and we will release our resources at https://github.com/VectorSpaceLab/OmniGen to foster future advancements.*
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import OmniGenTransformer2DModel
|
||||
|
||||
transformer = OmniGenTransformer2DModel.from_pretrained("Shitao/OmniGen-v1-diffusers", subfolder="transformer", torch_dtype=torch.bfloat16)
|
||||
```
|
||||
|
||||
## OmniGenTransformer2DModel
|
||||
|
||||
[[autodoc]] OmniGenTransformer2DModel
|
||||
30
docs/source/en/api/models/wan_transformer_3d.md
Normal file
30
docs/source/en/api/models/wan_transformer_3d.md
Normal file
@@ -0,0 +1,30 @@
|
||||
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License. -->
|
||||
|
||||
# WanTransformer3DModel
|
||||
|
||||
A Diffusion Transformer model for 3D video-like data was introduced in [Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
|
||||
|
||||
The model can be loaded with the following code snippet.
|
||||
|
||||
```python
|
||||
from diffusers import WanTransformer3DModel
|
||||
|
||||
transformer = WanTransformer3DModel.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", subfolder="transformer", torch_dtype=torch.bfloat16)
|
||||
```
|
||||
|
||||
## WanTransformer3DModel
|
||||
|
||||
[[autodoc]] WanTransformer3DModel
|
||||
|
||||
## Transformer2DModelOutput
|
||||
|
||||
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput
|
||||
@@ -29,3 +29,43 @@ Customized normalization layers for supporting various models in 🤗 Diffusers.
|
||||
## AdaGroupNorm
|
||||
|
||||
[[autodoc]] models.normalization.AdaGroupNorm
|
||||
|
||||
## AdaLayerNormContinuous
|
||||
|
||||
[[autodoc]] models.normalization.AdaLayerNormContinuous
|
||||
|
||||
## RMSNorm
|
||||
|
||||
[[autodoc]] models.normalization.RMSNorm
|
||||
|
||||
## GlobalResponseNorm
|
||||
|
||||
[[autodoc]] models.normalization.GlobalResponseNorm
|
||||
|
||||
|
||||
## LuminaLayerNormContinuous
|
||||
[[autodoc]] models.normalization.LuminaLayerNormContinuous
|
||||
|
||||
## SD35AdaLayerNormZeroX
|
||||
[[autodoc]] models.normalization.SD35AdaLayerNormZeroX
|
||||
|
||||
## AdaLayerNormZeroSingle
|
||||
[[autodoc]] models.normalization.AdaLayerNormZeroSingle
|
||||
|
||||
## LuminaRMSNormZero
|
||||
[[autodoc]] models.normalization.LuminaRMSNormZero
|
||||
|
||||
## LpNorm
|
||||
[[autodoc]] models.normalization.LpNorm
|
||||
|
||||
## CogView3PlusAdaLayerNormZeroTextImage
|
||||
[[autodoc]] models.normalization.CogView3PlusAdaLayerNormZeroTextImage
|
||||
|
||||
## CogVideoXLayerNormZero
|
||||
[[autodoc]] models.normalization.CogVideoXLayerNormZero
|
||||
|
||||
## MochiRMSNormZero
|
||||
[[autodoc]] models.transformers.transformer_mochi.MochiRMSNormZero
|
||||
|
||||
## MochiRMSNorm
|
||||
[[autodoc]] models.normalization.MochiRMSNorm
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Text-to-Video Generation with AnimateDiff
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
## Overview
|
||||
|
||||
[AnimateDiff: Animate Your Personalized Text-to-Image Diffusion Models without Specific Tuning](https://arxiv.org/abs/2307.04725) by Yuwei Guo, Ceyuan Yang, Anyi Rao, Yaohui Wang, Yu Qiao, Dahua Lin, Bo Dai.
|
||||
|
||||
@@ -89,6 +89,23 @@ image = pipeline(prompt).images[0]
|
||||
image.save("auraflow.png")
|
||||
```
|
||||
|
||||
## Support for `torch.compile()`
|
||||
|
||||
AuraFlow can be compiled with `torch.compile()` to speed up inference latency even for different resolutions. First, install PyTorch nightly following the instructions from [here](https://pytorch.org/). The snippet below shows the changes needed to enable this:
|
||||
|
||||
```diff
|
||||
+ torch.fx.experimental._config.use_duck_shape = False
|
||||
+ pipeline.transformer = torch.compile(
|
||||
pipeline.transformer, fullgraph=True, dynamic=True
|
||||
)
|
||||
```
|
||||
|
||||
Specifying `use_duck_shape` to be `False` instructs the compiler if it should use the same symbolic variable to represent input sizes that are the same. For more details, check out [this comment](https://github.com/huggingface/diffusers/pull/11327#discussion_r2047659790).
|
||||
|
||||
This enables from 100% (on low resolutions) to a 30% (on 1536x1536 resolution) speed improvements.
|
||||
|
||||
Thanks to [AstraliteHeart](https://github.com/huggingface/diffusers/pull/11297/) who helped us rewrite the [`AuraFlowTransformer2DModel`] class so that the above works for different resolutions ([PR](https://github.com/huggingface/diffusers/pull/11297/)).
|
||||
|
||||
## AuraFlowPipeline
|
||||
|
||||
[[autodoc]] AuraFlowPipeline
|
||||
|
||||
@@ -15,6 +15,10 @@
|
||||
|
||||
# CogVideoX
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer](https://arxiv.org/abs/2408.06072) from Tsinghua University & ZhipuAI, by Zhuoyi Yang, Jiayan Teng, Wendi Zheng, Ming Ding, Shiyu Huang, Jiazheng Xu, Yuanming Yang, Wenyi Hong, Xiaohan Zhang, Guanyu Feng, Da Yin, Xiaotao Gu, Yuxuan Zhang, Weihan Wang, Yean Cheng, Ting Liu, Bin Xu, Yuxiao Dong, Jie Tang.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
34
docs/source/en/api/pipelines/cogview4.md
Normal file
34
docs/source/en/api/pipelines/cogview4.md
Normal file
@@ -0,0 +1,34 @@
|
||||
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License.
|
||||
-->
|
||||
|
||||
# CogView4
|
||||
|
||||
<Tip>
|
||||
|
||||
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
|
||||
|
||||
</Tip>
|
||||
|
||||
This pipeline was contributed by [zRzRzRzRzRzRzR](https://github.com/zRzRzRzRzRzRzR). The original codebase can be found [here](https://huggingface.co/THUDM). The original weights can be found under [hf.co/THUDM](https://huggingface.co/THUDM).
|
||||
|
||||
## CogView4Pipeline
|
||||
|
||||
[[autodoc]] CogView4Pipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## CogView4PipelineOutput
|
||||
|
||||
[[autodoc]] pipelines.cogview4.pipeline_output.CogView4PipelineOutput
|
||||
@@ -15,6 +15,10 @@
|
||||
|
||||
# ConsisID
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[Identity-Preserving Text-to-Video Generation by Frequency Decomposition](https://arxiv.org/abs/2411.17440) from Peking University & University of Rochester & etc, by Shenghai Yuan, Jinfa Huang, Xianyi He, Yunyang Ge, Yujun Shi, Liuhan Chen, Jiebo Luo, Li Yuan.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# FluxControlInpaint
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
FluxControlInpaintPipeline is an implementation of Inpainting for Flux.1 Depth/Canny models. It is a pipeline that allows you to inpaint images using the Flux.1 Depth/Canny models. The pipeline takes an image and a mask as input and returns the inpainted image.
|
||||
|
||||
FLUX.1 Depth and Canny [dev] is a 12 billion parameter rectified flow transformer capable of generating an image based on a text description while following the structure of a given input image. **This is not a ControlNet model**.
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# ControlNet
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
|
||||
|
||||
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# ControlNet with Flux.1
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
FluxControlNetPipeline is an implementation of ControlNet for Flux.1.
|
||||
|
||||
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
|
||||
|
||||
36
docs/source/en/api/pipelines/controlnet_sana.md
Normal file
36
docs/source/en/api/pipelines/controlnet_sana.md
Normal file
@@ -0,0 +1,36 @@
|
||||
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# ControlNet
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
|
||||
|
||||
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
|
||||
|
||||
This pipeline was contributed by [ishan24](https://huggingface.co/ishan24). ❤️
|
||||
The original codebase can be found at [NVlabs/Sana](https://github.com/NVlabs/Sana), and you can find official ControlNet checkpoints on [Efficient-Large-Model's](https://huggingface.co/Efficient-Large-Model) Hub profile.
|
||||
|
||||
## SanaControlNetPipeline
|
||||
[[autodoc]] SanaControlNetPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## SanaPipelineOutput
|
||||
[[autodoc]] pipelines.sana.pipeline_output.SanaPipelineOutput
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# ControlNet with Stable Diffusion 3
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
StableDiffusion3ControlNetPipeline is an implementation of ControlNet for Stable Diffusion 3.
|
||||
|
||||
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# ControlNet with Stable Diffusion XL
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
|
||||
|
||||
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# ControlNetUnion
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
ControlNetUnionModel is an implementation of ControlNet for Stable Diffusion XL.
|
||||
|
||||
The ControlNet model was introduced in [ControlNetPlus](https://github.com/xinsir6/ControlNetPlus) by xinsir6. It supports multiple conditioning inputs without increasing computation.
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# ControlNet-XS
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
|
||||
|
||||
Like the original ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
|
||||
|
||||
@@ -12,6 +12,11 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# DeepFloyd IF
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
|
||||
</div>
|
||||
|
||||
## Overview
|
||||
|
||||
DeepFloyd IF is a novel state-of-the-art open-source text-to-image model with a high degree of photorealism and language understanding.
|
||||
|
||||
88
docs/source/en/api/pipelines/easyanimate.md
Normal file
88
docs/source/en/api/pipelines/easyanimate.md
Normal file
@@ -0,0 +1,88 @@
|
||||
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License.
|
||||
-->
|
||||
|
||||
# EasyAnimate
|
||||
[EasyAnimate](https://github.com/aigc-apps/EasyAnimate) by Alibaba PAI.
|
||||
|
||||
The description from it's GitHub page:
|
||||
*EasyAnimate is a pipeline based on the transformer architecture, designed for generating AI images and videos, and for training baseline models and Lora models for Diffusion Transformer. We support direct prediction from pre-trained EasyAnimate models, allowing for the generation of videos with various resolutions, approximately 6 seconds in length, at 8fps (EasyAnimateV5.1, 1 to 49 frames). Additionally, users can train their own baseline and Lora models for specific style transformations.*
|
||||
|
||||
This pipeline was contributed by [bubbliiiing](https://github.com/bubbliiiing). The original codebase can be found [here](https://huggingface.co/alibaba-pai). The original weights can be found under [hf.co/alibaba-pai](https://huggingface.co/alibaba-pai).
|
||||
|
||||
There are two official EasyAnimate checkpoints for text-to-video and video-to-video.
|
||||
|
||||
| checkpoints | recommended inference dtype |
|
||||
|:---:|:---:|
|
||||
| [`alibaba-pai/EasyAnimateV5.1-12b-zh`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh) | torch.float16 |
|
||||
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-InP`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-InP) | torch.float16 |
|
||||
|
||||
There is one official EasyAnimate checkpoints available for image-to-video and video-to-video.
|
||||
|
||||
| checkpoints | recommended inference dtype |
|
||||
|:---:|:---:|
|
||||
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-InP`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-InP) | torch.float16 |
|
||||
|
||||
There are two official EasyAnimate checkpoints available for control-to-video.
|
||||
|
||||
| checkpoints | recommended inference dtype |
|
||||
|:---:|:---:|
|
||||
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-Control`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-Control) | torch.float16 |
|
||||
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-Control-Camera`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-Control-Camera) | torch.float16 |
|
||||
|
||||
For the EasyAnimateV5.1 series:
|
||||
- Text-to-video (T2V) and Image-to-video (I2V) works for multiple resolutions. The width and height can vary from 256 to 1024.
|
||||
- Both T2V and I2V models support generation with 1~49 frames and work best at this value. Exporting videos at 8 FPS is recommended.
|
||||
|
||||
## Quantization
|
||||
|
||||
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
|
||||
|
||||
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`EasyAnimatePipeline`] for inference with bitsandbytes.
|
||||
|
||||
```py
|
||||
import torch
|
||||
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, EasyAnimateTransformer3DModel, EasyAnimatePipeline
|
||||
from diffusers.utils import export_to_video
|
||||
|
||||
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True)
|
||||
transformer_8bit = EasyAnimateTransformer3DModel.from_pretrained(
|
||||
"alibaba-pai/EasyAnimateV5.1-12b-zh",
|
||||
subfolder="transformer",
|
||||
quantization_config=quant_config,
|
||||
torch_dtype=torch.float16,
|
||||
)
|
||||
|
||||
pipeline = EasyAnimatePipeline.from_pretrained(
|
||||
"alibaba-pai/EasyAnimateV5.1-12b-zh",
|
||||
transformer=transformer_8bit,
|
||||
torch_dtype=torch.float16,
|
||||
device_map="balanced",
|
||||
)
|
||||
|
||||
prompt = "A cat walks on the grass, realistic style."
|
||||
negative_prompt = "bad detailed"
|
||||
video = pipeline(prompt=prompt, negative_prompt=negative_prompt, num_frames=49, num_inference_steps=30).frames[0]
|
||||
export_to_video(video, "cat.mp4", fps=8)
|
||||
```
|
||||
|
||||
## EasyAnimatePipeline
|
||||
|
||||
[[autodoc]] EasyAnimatePipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## EasyAnimatePipelineOutput
|
||||
|
||||
[[autodoc]] pipelines.easyanimate.pipeline_output.EasyAnimatePipelineOutput
|
||||
@@ -12,6 +12,11 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Flux
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
|
||||
</div>
|
||||
|
||||
Flux is a series of text-to-image generation models based on diffusion transformers. To know more about Flux, check out the original [blog post](https://blackforestlabs.ai/announcing-black-forest-labs/) by the creators of Flux, Black Forest Labs.
|
||||
|
||||
Original model checkpoints for Flux can be found [here](https://huggingface.co/black-forest-labs). Original inference code can be found [here](https://github.com/black-forest-labs/flux).
|
||||
@@ -309,7 +314,120 @@ image.save("output.png")
|
||||
|
||||
When unloading the Control LoRA weights, call `pipe.unload_lora_weights(reset_to_overwritten_params=True)` to reset the `pipe.transformer` completely back to its original form. The resultant pipeline can then be used with methods like [`DiffusionPipeline.from_pipe`]. More details about this argument are available in [this PR](https://github.com/huggingface/diffusers/pull/10397).
|
||||
|
||||
## Running FP16 inference
|
||||
## IP-Adapter
|
||||
|
||||
<Tip>
|
||||
|
||||
Check out [IP-Adapter](../../../using-diffusers/ip_adapter) to learn more about how IP-Adapters work.
|
||||
|
||||
</Tip>
|
||||
|
||||
An IP-Adapter lets you prompt Flux with images, in addition to the text prompt. This is especially useful when describing complex concepts that are difficult to articulate through text alone and you have reference images.
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import FluxPipeline
|
||||
from diffusers.utils import load_image
|
||||
|
||||
pipe = FluxPipeline.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-dev", torch_dtype=torch.bfloat16
|
||||
).to("cuda")
|
||||
|
||||
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flux_ip_adapter_input.jpg").resize((1024, 1024))
|
||||
|
||||
pipe.load_ip_adapter(
|
||||
"XLabs-AI/flux-ip-adapter",
|
||||
weight_name="ip_adapter.safetensors",
|
||||
image_encoder_pretrained_model_name_or_path="openai/clip-vit-large-patch14"
|
||||
)
|
||||
pipe.set_ip_adapter_scale(1.0)
|
||||
|
||||
image = pipe(
|
||||
width=1024,
|
||||
height=1024,
|
||||
prompt="wearing sunglasses",
|
||||
negative_prompt="",
|
||||
true_cfg=4.0,
|
||||
generator=torch.Generator().manual_seed(4444),
|
||||
ip_adapter_image=image,
|
||||
).images[0]
|
||||
|
||||
image.save('flux_ip_adapter_output.jpg')
|
||||
```
|
||||
|
||||
<div class="justify-center">
|
||||
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flux_ip_adapter_output.jpg"/>
|
||||
<figcaption class="mt-2 text-sm text-center text-gray-500">IP-Adapter examples with prompt "wearing sunglasses"</figcaption>
|
||||
</div>
|
||||
|
||||
## Optimize
|
||||
|
||||
Flux is a very large model and requires ~50GB of RAM/VRAM to load all the modeling components. Enable some of the optimizations below to lower the memory requirements.
|
||||
|
||||
### Group offloading
|
||||
|
||||
[Group offloading](../../optimization/memory#group-offloading) lowers VRAM usage by offloading groups of internal layers rather than the whole model or weights. You need to use [`~hooks.apply_group_offloading`] on all the model components of a pipeline. The `offload_type` parameter allows you to toggle between block and leaf-level offloading. Setting it to `leaf_level` offloads the lowest leaf-level parameters to the CPU instead of offloading at the module-level.
|
||||
|
||||
On CUDA devices that support asynchronous data streaming, set `use_stream=True` to overlap data transfer and computation to accelerate inference.
|
||||
|
||||
> [!TIP]
|
||||
> It is possible to mix block and leaf-level offloading for different components in a pipeline.
|
||||
|
||||
```py
|
||||
import torch
|
||||
from diffusers import FluxPipeline
|
||||
from diffusers.hooks import apply_group_offloading
|
||||
|
||||
model_id = "black-forest-labs/FLUX.1-dev"
|
||||
dtype = torch.bfloat16
|
||||
pipe = FluxPipeline.from_pretrained(
|
||||
model_id,
|
||||
torch_dtype=dtype,
|
||||
)
|
||||
|
||||
apply_group_offloading(
|
||||
pipe.transformer,
|
||||
offload_type="leaf_level",
|
||||
offload_device=torch.device("cpu"),
|
||||
onload_device=torch.device("cuda"),
|
||||
use_stream=True,
|
||||
)
|
||||
apply_group_offloading(
|
||||
pipe.text_encoder,
|
||||
offload_device=torch.device("cpu"),
|
||||
onload_device=torch.device("cuda"),
|
||||
offload_type="leaf_level",
|
||||
use_stream=True,
|
||||
)
|
||||
apply_group_offloading(
|
||||
pipe.text_encoder_2,
|
||||
offload_device=torch.device("cpu"),
|
||||
onload_device=torch.device("cuda"),
|
||||
offload_type="leaf_level",
|
||||
use_stream=True,
|
||||
)
|
||||
apply_group_offloading(
|
||||
pipe.vae,
|
||||
offload_device=torch.device("cpu"),
|
||||
onload_device=torch.device("cuda"),
|
||||
offload_type="leaf_level",
|
||||
use_stream=True,
|
||||
)
|
||||
|
||||
prompt="A cat wearing sunglasses and working as a lifeguard at pool."
|
||||
|
||||
generator = torch.Generator().manual_seed(181201)
|
||||
image = pipe(
|
||||
prompt,
|
||||
width=576,
|
||||
height=1024,
|
||||
num_inference_steps=30,
|
||||
generator=generator
|
||||
).images[0]
|
||||
image
|
||||
```
|
||||
|
||||
### Running FP16 inference
|
||||
|
||||
Flux can generate high-quality images with FP16 (i.e. to accelerate inference on Turing/Volta GPUs) but produces different outputs compared to FP32/BF16. The issue is that some activations in the text encoders have to be clipped when running in FP16, which affects the overall image. Forcing text encoders to run with FP32 inference thus removes this output difference. See [here](https://github.com/huggingface/diffusers/pull/9097#issuecomment-2272292516) for details.
|
||||
|
||||
@@ -338,7 +456,7 @@ out = pipe(
|
||||
out.save("image.png")
|
||||
```
|
||||
|
||||
## Quantization
|
||||
### Quantization
|
||||
|
||||
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
|
||||
|
||||
|
||||
43
docs/source/en/api/pipelines/hidream.md
Normal file
43
docs/source/en/api/pipelines/hidream.md
Normal file
@@ -0,0 +1,43 @@
|
||||
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License. -->
|
||||
|
||||
# HiDreamImage
|
||||
|
||||
[HiDream-I1](https://huggingface.co/HiDream-ai) by HiDream.ai
|
||||
|
||||
<Tip>
|
||||
|
||||
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
|
||||
|
||||
</Tip>
|
||||
|
||||
## Available models
|
||||
|
||||
The following models are available for the [`HiDreamImagePipeline`](text-to-image) pipeline:
|
||||
|
||||
| Model name | Description |
|
||||
|:---|:---|
|
||||
| [`HiDream-ai/HiDream-I1-Full`](https://huggingface.co/HiDream-ai/HiDream-I1-Full) | - |
|
||||
| [`HiDream-ai/HiDream-I1-Dev`](https://huggingface.co/HiDream-ai/HiDream-I1-Dev) | - |
|
||||
| [`HiDream-ai/HiDream-I1-Fast`](https://huggingface.co/HiDream-ai/HiDream-I1-Fast) | - |
|
||||
|
||||
## HiDreamImagePipeline
|
||||
|
||||
[[autodoc]] HiDreamImagePipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## HiDreamImagePipelineOutput
|
||||
|
||||
[[autodoc]] pipelines.hidream_image.pipeline_output.HiDreamImagePipelineOutput
|
||||
@@ -14,6 +14,10 @@
|
||||
|
||||
# HunyuanVideo
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[HunyuanVideo](https://www.arxiv.org/abs/2412.03603) by Tencent.
|
||||
|
||||
*Recent advancements in video generation have significantly impacted daily life for both individuals and industries. However, the leading video generation models remain closed-source, resulting in a notable performance gap between industry capabilities and those available to the public. In this report, we introduce HunyuanVideo, an innovative open-source video foundation model that demonstrates performance in video generation comparable to, or even surpassing, that of leading closed-source models. HunyuanVideo encompasses a comprehensive framework that integrates several key elements, including data curation, advanced architectural design, progressive model scaling and training, and an efficient infrastructure tailored for large-scale model training and inference. As a result, we successfully trained a video generative model with over 13 billion parameters, making it the largest among all open-source models. We conducted extensive experiments and implemented a series of targeted designs to ensure high visual quality, motion dynamics, text-video alignment, and advanced filming techniques. According to evaluations by professionals, HunyuanVideo outperforms previous state-of-the-art models, including Runway Gen-3, Luma 1.6, and three top-performing Chinese video generative models. By releasing the code for the foundation model and its applications, we aim to bridge the gap between closed-source and open-source communities. This initiative will empower individuals within the community to experiment with their ideas, fostering a more dynamic and vibrant video generation ecosystem. The code is publicly available at [this https URL](https://github.com/tencent/HunyuanVideo).*
|
||||
@@ -32,6 +36,23 @@ Recommendations for inference:
|
||||
- For smaller resolution videos, try lower values of `shift` (between `2.0` to `5.0`) in the [Scheduler](https://huggingface.co/docs/diffusers/main/en/api/schedulers/flow_match_euler_discrete#diffusers.FlowMatchEulerDiscreteScheduler.shift). For larger resolution images, try higher values (between `7.0` and `12.0`). The default value is `7.0` for HunyuanVideo.
|
||||
- For more information about supported resolutions and other details, please refer to the original repository [here](https://github.com/Tencent/HunyuanVideo/).
|
||||
|
||||
## Available models
|
||||
|
||||
The following models are available for the [`HunyuanVideoPipeline`](text-to-video) pipeline:
|
||||
|
||||
| Model name | Description |
|
||||
|:---|:---|
|
||||
| [`hunyuanvideo-community/HunyuanVideo`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo) | Official HunyuanVideo (guidance-distilled). Performs best at multiple resolutions and frames. Performs best with `guidance_scale=6.0`, `true_cfg_scale=1.0` and without a negative prompt. |
|
||||
| [`https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-T2V`](https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-T2V) | Skywork's custom finetune of HunyuanVideo (de-distilled). Performs best with `97x544x960` resolution, `guidance_scale=1.0`, `true_cfg_scale=6.0` and a negative prompt. |
|
||||
|
||||
The following models are available for the image-to-video pipeline:
|
||||
|
||||
| Model name | Description |
|
||||
|:---|:---|
|
||||
| [`Skywork/SkyReels-V1-Hunyuan-I2V`](https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-I2V) | Skywork's custom finetune of HunyuanVideo (de-distilled). Performs best with `97x544x960` resolution. Performs best at `97x544x960` resolution, `guidance_scale=1.0`, `true_cfg_scale=6.0` and a negative prompt. |
|
||||
| [`hunyuanvideo-community/HunyuanVideo-I2V-33ch`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo-I2V) | Tecent's official HunyuanVideo 33-channel I2V model. Performs best at resolutions of 480, 720, 960, 1280. A higher `shift` value when initializing the scheduler is recommended (good values are between 7 and 20). |
|
||||
| [`hunyuanvideo-community/HunyuanVideo-I2V`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo-I2V) | Tecent's official HunyuanVideo 16-channel I2V model. Performs best at resolutions of 480, 720, 960, 1280. A higher `shift` value when initializing the scheduler is recommended (good values are between 7 and 20) |
|
||||
|
||||
## Quantization
|
||||
|
||||
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
|
||||
|
||||
@@ -9,6 +9,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Kandinsky 3
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
Kandinsky 3 is created by [Vladimir Arkhipkin](https://github.com/oriBetelgeuse),[Anastasia Maltseva](https://github.com/NastyaMittseva),[Igor Pavlov](https://github.com/boomb0om),[Andrei Filatov](https://github.com/anvilarth),[Arseniy Shakhmatov](https://github.com/cene555),[Andrey Kuznetsov](https://github.com/kuznetsoffandrey),[Denis Dimitrov](https://github.com/denndimitrov), [Zein Shaheen](https://github.com/zeinsh)
|
||||
|
||||
The description from it's GitHub page:
|
||||
|
||||
@@ -12,6 +12,11 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Kolors: Effective Training of Diffusion Model for Photorealistic Text-to-Image Synthesis
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
|
||||
</div>
|
||||
|
||||

|
||||
|
||||
Kolors is a large-scale text-to-image generation model based on latent diffusion, developed by [the Kuaishou Kolors team](https://github.com/Kwai-Kolors/Kolors). Trained on billions of text-image pairs, Kolors exhibits significant advantages over both open-source and closed-source models in visual quality, complex semantic accuracy, and text rendering for both Chinese and English characters. Furthermore, Kolors supports both Chinese and English inputs, demonstrating strong performance in understanding and generating Chinese-specific content. For more details, please refer to this [technical report](https://github.com/Kwai-Kolors/Kolors/blob/master/imgs/Kolors_paper.pdf).
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Latent Consistency Models
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
Latent Consistency Models (LCMs) were proposed in [Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference](https://huggingface.co/papers/2310.04378) by Simian Luo, Yiqin Tan, Longbo Huang, Jian Li, and Hang Zhao.
|
||||
|
||||
The abstract of the paper is as follows:
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# LEDITS++
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
LEDITS++ was proposed in [LEDITS++: Limitless Image Editing using Text-to-Image Models](https://huggingface.co/papers/2311.16711) by Manuel Brack, Felix Friedrich, Katharina Kornmeier, Linoy Tsaban, Patrick Schramowski, Kristian Kersting, Apolinário Passos.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
@@ -14,6 +14,11 @@
|
||||
|
||||
# LTX Video
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
|
||||
</div>
|
||||
|
||||
[LTX Video](https://huggingface.co/Lightricks/LTX-Video) is the first DiT-based video generation model capable of generating high-quality videos in real-time. It produces 24 FPS videos at a 768x512 resolution faster than they can be watched. Trained on a large-scale dataset of diverse videos, the model generates high-resolution videos with realistic and varied content. We provide a model for both text-to-video as well as image + text-to-video usecases.
|
||||
|
||||
<Tip>
|
||||
@@ -28,6 +33,7 @@ Available models:
|
||||
|:-------------:|:-----------------:|
|
||||
| [`LTX Video 0.9.0`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.safetensors) | `torch.bfloat16` |
|
||||
| [`LTX Video 0.9.1`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.1.safetensors) | `torch.bfloat16` |
|
||||
| [`LTX Video 0.9.5`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.5.safetensors) | `torch.bfloat16` |
|
||||
|
||||
Note: The recommended dtype is for the transformer component. The VAE and text encoders can be either `torch.float32`, `torch.bfloat16` or `torch.float16` but the recommended dtype is `torch.bfloat16` as used in the original repository.
|
||||
|
||||
@@ -192,6 +198,12 @@ export_to_video(video, "ship.mp4", fps=24)
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## LTXConditionPipeline
|
||||
|
||||
[[autodoc]] LTXConditionPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## LTXPipelineOutput
|
||||
|
||||
[[autodoc]] pipelines.ltx.pipeline_output.LTXPipelineOutput
|
||||
|
||||
@@ -58,10 +58,10 @@ Use [`torch.compile`](https://huggingface.co/docs/diffusers/main/en/tutorials/fa
|
||||
First, load the pipeline:
|
||||
|
||||
```python
|
||||
from diffusers import LuminaText2ImgPipeline
|
||||
from diffusers import LuminaPipeline
|
||||
import torch
|
||||
|
||||
pipeline = LuminaText2ImgPipeline.from_pretrained(
|
||||
pipeline = LuminaPipeline.from_pretrained(
|
||||
"Alpha-VLLM/Lumina-Next-SFT-diffusers", torch_dtype=torch.bfloat16
|
||||
).to("cuda")
|
||||
```
|
||||
@@ -86,11 +86,11 @@ image = pipeline(prompt="Upper body of a young woman in a Victorian-era outfit w
|
||||
|
||||
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
|
||||
|
||||
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`LuminaText2ImgPipeline`] for inference with bitsandbytes.
|
||||
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`LuminaPipeline`] for inference with bitsandbytes.
|
||||
|
||||
```py
|
||||
import torch
|
||||
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, Transformer2DModel, LuminaText2ImgPipeline
|
||||
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, Transformer2DModel, LuminaPipeline
|
||||
from transformers import BitsAndBytesConfig as BitsAndBytesConfig, T5EncoderModel
|
||||
|
||||
quant_config = BitsAndBytesConfig(load_in_8bit=True)
|
||||
@@ -109,7 +109,7 @@ transformer_8bit = Transformer2DModel.from_pretrained(
|
||||
torch_dtype=torch.float16,
|
||||
)
|
||||
|
||||
pipeline = LuminaText2ImgPipeline.from_pretrained(
|
||||
pipeline = LuminaPipeline.from_pretrained(
|
||||
"Alpha-VLLM/Lumina-Next-SFT-diffusers",
|
||||
text_encoder=text_encoder_8bit,
|
||||
transformer=transformer_8bit,
|
||||
@@ -122,9 +122,9 @@ image = pipeline(prompt).images[0]
|
||||
image.save("lumina.png")
|
||||
```
|
||||
|
||||
## LuminaText2ImgPipeline
|
||||
## LuminaPipeline
|
||||
|
||||
[[autodoc]] LuminaText2ImgPipeline
|
||||
[[autodoc]] LuminaPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
|
||||
87
docs/source/en/api/pipelines/lumina2.md
Normal file
87
docs/source/en/api/pipelines/lumina2.md
Normal file
@@ -0,0 +1,87 @@
|
||||
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License. -->
|
||||
|
||||
# Lumina2
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[Lumina Image 2.0: A Unified and Efficient Image Generative Model](https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0) is a 2 billion parameter flow-based diffusion transformer capable of generating diverse images from text descriptions.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
*We introduce Lumina-Image 2.0, an advanced text-to-image model that surpasses previous state-of-the-art methods across multiple benchmarks, while also shedding light on its potential to evolve into a generalist vision intelligence model. Lumina-Image 2.0 exhibits three key properties: (1) Unification – it adopts a unified architecture that treats text and image tokens as a joint sequence, enabling natural cross-modal interactions and facilitating task expansion. Besides, since high-quality captioners can provide semantically better-aligned text-image training pairs, we introduce a unified captioning system, UniCaptioner, which generates comprehensive and precise captions for the model. This not only accelerates model convergence but also enhances prompt adherence, variable-length prompt handling, and task generalization via prompt templates. (2) Efficiency – to improve the efficiency of the unified architecture, we develop a set of optimization techniques that improve semantic learning and fine-grained texture generation during training while incorporating inference-time acceleration strategies without compromising image quality. (3) Transparency – we open-source all training details, code, and models to ensure full reproducibility, aiming to bridge the gap between well-resourced closed-source research teams and independent developers.*
|
||||
|
||||
<Tip>
|
||||
|
||||
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
|
||||
|
||||
</Tip>
|
||||
|
||||
## Using Single File loading with Lumina Image 2.0
|
||||
|
||||
Single file loading for Lumina Image 2.0 is available for the `Lumina2Transformer2DModel`
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import Lumina2Transformer2DModel, Lumina2Pipeline
|
||||
|
||||
ckpt_path = "https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0/blob/main/consolidated.00-of-01.pth"
|
||||
transformer = Lumina2Transformer2DModel.from_single_file(
|
||||
ckpt_path, torch_dtype=torch.bfloat16
|
||||
)
|
||||
|
||||
pipe = Lumina2Pipeline.from_pretrained(
|
||||
"Alpha-VLLM/Lumina-Image-2.0", transformer=transformer, torch_dtype=torch.bfloat16
|
||||
)
|
||||
pipe.enable_model_cpu_offload()
|
||||
image = pipe(
|
||||
"a cat holding a sign that says hello",
|
||||
generator=torch.Generator("cpu").manual_seed(0),
|
||||
).images[0]
|
||||
image.save("lumina-single-file.png")
|
||||
|
||||
```
|
||||
|
||||
## Using GGUF Quantized Checkpoints with Lumina Image 2.0
|
||||
|
||||
GGUF Quantized checkpoints for the `Lumina2Transformer2DModel` can be loaded via `from_single_file` with the `GGUFQuantizationConfig`
|
||||
|
||||
```python
|
||||
from diffusers import Lumina2Transformer2DModel, Lumina2Pipeline, GGUFQuantizationConfig
|
||||
|
||||
ckpt_path = "https://huggingface.co/calcuis/lumina-gguf/blob/main/lumina2-q4_0.gguf"
|
||||
transformer = Lumina2Transformer2DModel.from_single_file(
|
||||
ckpt_path,
|
||||
quantization_config=GGUFQuantizationConfig(compute_dtype=torch.bfloat16),
|
||||
torch_dtype=torch.bfloat16,
|
||||
)
|
||||
|
||||
pipe = Lumina2Pipeline.from_pretrained(
|
||||
"Alpha-VLLM/Lumina-Image-2.0", transformer=transformer, torch_dtype=torch.bfloat16
|
||||
)
|
||||
pipe.enable_model_cpu_offload()
|
||||
image = pipe(
|
||||
"a cat holding a sign that says hello",
|
||||
generator=torch.Generator("cpu").manual_seed(0),
|
||||
).images[0]
|
||||
image.save("lumina-gguf.png")
|
||||
```
|
||||
|
||||
## Lumina2Pipeline
|
||||
|
||||
[[autodoc]] Lumina2Pipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -1,4 +1,6 @@
|
||||
<!--Copyright 2024 Marigold authors and The HuggingFace Team. All rights reserved.
|
||||
<!--
|
||||
Copyright 2023-2025 Marigold Team, ETH Zürich. All rights reserved.
|
||||
Copyright 2024-2025 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
@@ -10,67 +12,120 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Marigold Pipelines for Computer Vision Tasks
|
||||
# Marigold Computer Vision
|
||||
|
||||

|
||||
|
||||
Marigold was proposed in [Repurposing Diffusion-Based Image Generators for Monocular Depth Estimation](https://huggingface.co/papers/2312.02145), a CVPR 2024 Oral paper by [Bingxin Ke](http://www.kebingxin.com/), [Anton Obukhov](https://www.obukhov.ai/), [Shengyu Huang](https://shengyuh.github.io/), [Nando Metzger](https://nandometzger.github.io/), [Rodrigo Caye Daudt](https://rcdaudt.github.io/), and [Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
|
||||
The idea is to repurpose the rich generative prior of Text-to-Image Latent Diffusion Models (LDMs) for traditional computer vision tasks.
|
||||
Initially, this idea was explored to fine-tune Stable Diffusion for Monocular Depth Estimation, as shown in the teaser above.
|
||||
Later,
|
||||
- [Tianfu Wang](https://tianfwang.github.io/) trained the first Latent Consistency Model (LCM) of Marigold, which unlocked fast single-step inference;
|
||||
- [Kevin Qu](https://www.linkedin.com/in/kevin-qu-b3417621b/?locale=en_US) extended the approach to Surface Normals Estimation;
|
||||
- [Anton Obukhov](https://www.obukhov.ai/) contributed the pipelines and documentation into diffusers (enabled and supported by [YiYi Xu](https://yiyixuxu.github.io/) and [Sayak Paul](https://sayak.dev/)).
|
||||
Marigold was proposed in
|
||||
[Repurposing Diffusion-Based Image Generators for Monocular Depth Estimation](https://huggingface.co/papers/2312.02145),
|
||||
a CVPR 2024 Oral paper by
|
||||
[Bingxin Ke](http://www.kebingxin.com/),
|
||||
[Anton Obukhov](https://www.obukhov.ai/),
|
||||
[Shengyu Huang](https://shengyuh.github.io/),
|
||||
[Nando Metzger](https://nandometzger.github.io/),
|
||||
[Rodrigo Caye Daudt](https://rcdaudt.github.io/), and
|
||||
[Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
|
||||
The core idea is to **repurpose the generative prior of Text-to-Image Latent Diffusion Models (LDMs) for traditional
|
||||
computer vision tasks**.
|
||||
This approach was explored by fine-tuning Stable Diffusion for **Monocular Depth Estimation**, as demonstrated in the
|
||||
teaser above.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
*Monocular depth estimation is a fundamental computer vision task. Recovering 3D depth from a single image is geometrically ill-posed and requires scene understanding, so it is not surprising that the rise of deep learning has led to a breakthrough. The impressive progress of monocular depth estimators has mirrored the growth in model capacity, from relatively modest CNNs to large Transformer architectures. Still, monocular depth estimators tend to struggle when presented with images with unfamiliar content and layout, since their knowledge of the visual world is restricted by the data seen during training, and challenged by zero-shot generalization to new domains. This motivates us to explore whether the extensive priors captured in recent generative diffusion models can enable better, more generalizable depth estimation. We introduce Marigold, a method for affine-invariant monocular depth estimation that is derived from Stable Diffusion and retains its rich prior knowledge. The estimator can be fine-tuned in a couple of days on a single GPU using only synthetic training data. It delivers state-of-the-art performance across a wide range of datasets, including over 20% performance gains in specific cases. Project page: https://marigoldmonodepth.github.io.*
|
||||
|
||||
## Available Pipelines
|
||||
|
||||
Each pipeline supports one Computer Vision task, which takes an input RGB image as input and produces a *prediction* of the modality of interest, such as a depth map of the input image.
|
||||
Currently, the following tasks are implemented:
|
||||
|
||||
| Pipeline | Predicted Modalities | Demos |
|
||||
|---------------------------------------------------------------------------------------------------------------------------------------------|------------------------------------------------------------------------------------------------------------------|:--------------------------------------------------------------------------------------------------------------------------------------------------:|
|
||||
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-lcm), [Slow Original Demo (DDIM)](https://huggingface.co/spaces/prs-eth/marigold) |
|
||||
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-normals-lcm) |
|
||||
|
||||
|
||||
## Available Checkpoints
|
||||
|
||||
The original checkpoints can be found under the [PRS-ETH](https://huggingface.co/prs-eth/) Hugging Face organization.
|
||||
Marigold was later extended in the follow-up paper,
|
||||
[Marigold: Affordable Adaptation of Diffusion-Based Image Generators for Image Analysis](https://huggingface.co/papers/2312.02145),
|
||||
authored by
|
||||
[Bingxin Ke](http://www.kebingxin.com/),
|
||||
[Kevin Qu](https://www.linkedin.com/in/kevin-qu-b3417621b/?locale=en_US),
|
||||
[Tianfu Wang](https://tianfwang.github.io/),
|
||||
[Nando Metzger](https://nandometzger.github.io/),
|
||||
[Shengyu Huang](https://shengyuh.github.io/),
|
||||
[Bo Li](https://www.linkedin.com/in/bobboli0202/),
|
||||
[Anton Obukhov](https://www.obukhov.ai/), and
|
||||
[Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
|
||||
This work expanded Marigold to support new modalities such as **Surface Normals** and **Intrinsic Image Decomposition**
|
||||
(IID), introduced a training protocol for **Latent Consistency Models** (LCM), and demonstrated **High-Resolution** (HR)
|
||||
processing capability.
|
||||
|
||||
<Tip>
|
||||
|
||||
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines. Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section [here](../../using-diffusers/svd#reduce-memory-usage).
|
||||
The early Marigold models (`v1-0` and earlier) were optimized for best results with at least 10 inference steps.
|
||||
LCM models were later developed to enable high-quality inference in just 1 to 4 steps.
|
||||
Marigold models `v1-1` and later use the DDIM scheduler to achieve optimal
|
||||
results in as few as 1 to 4 steps.
|
||||
|
||||
</Tip>
|
||||
|
||||
## Available Pipelines
|
||||
|
||||
Each pipeline is tailored for a specific computer vision task, processing an input RGB image and generating a
|
||||
corresponding prediction.
|
||||
Currently, the following computer vision tasks are implemented:
|
||||
|
||||
| Pipeline | Recommended Model Checkpoints | Spaces (Interactive Apps) | Predicted Modalities |
|
||||
|---------------------------------------------------------------------------------------------------------------------------------------------------|---------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|:------------------------------------------------------------------------------------:|------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
|
||||
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [prs-eth/marigold-depth-v1-1](https://huggingface.co/prs-eth/marigold-depth-v1-1) | [Depth Estimation](https://huggingface.co/spaces/prs-eth/marigold) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) |
|
||||
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [prs-eth/marigold-normals-v1-1](https://huggingface.co/prs-eth/marigold-normals-v1-1) | [Surface Normals Estimation](https://huggingface.co/spaces/prs-eth/marigold-normals) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) |
|
||||
| [MarigoldIntrinsicsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_intrinsics.py) | [prs-eth/marigold-iid-appearance-v1-1](https://huggingface.co/prs-eth/marigold-iid-appearance-v1-1),<br>[prs-eth/marigold-iid-lighting-v1-1](https://huggingface.co/prs-eth/marigold-iid-lighting-v1-1) | [Intrinsic Image Decomposition](https://huggingface.co/spaces/prs-eth/marigold-iid) | [Albedo](https://en.wikipedia.org/wiki/Albedo), [Materials](https://www.n.aiq3d.com/wiki/roughnessmetalnessao-map), [Lighting](https://en.wikipedia.org/wiki/Diffuse_reflection) |
|
||||
|
||||
## Available Checkpoints
|
||||
|
||||
All original checkpoints are available under the [PRS-ETH](https://huggingface.co/prs-eth/) organization on Hugging Face.
|
||||
They are designed for use with diffusers pipelines and the [original codebase](https://github.com/prs-eth/marigold), which can also be used to train
|
||||
new model checkpoints.
|
||||
The following is a summary of the recommended checkpoints, all of which produce reliable results with 1 to 4 steps.
|
||||
|
||||
| Checkpoint | Modality | Comment |
|
||||
|-----------------------------------------------------------------------------------------------------|--------------|--------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
|
||||
| [prs-eth/marigold-depth-v1-1](https://huggingface.co/prs-eth/marigold-depth-v1-1) | Depth | Affine-invariant depth prediction assigns each pixel a value between 0 (near plane) and 1 (far plane), with both planes determined by the model during inference. |
|
||||
| [prs-eth/marigold-normals-v0-1](https://huggingface.co/prs-eth/marigold-normals-v0-1) | Normals | The surface normals predictions are unit-length 3D vectors in the screen space camera, with values in the range from -1 to 1. |
|
||||
| [prs-eth/marigold-iid-appearance-v1-1](https://huggingface.co/prs-eth/marigold-iid-appearance-v1-1) | Intrinsics | InteriorVerse decomposition is comprised of Albedo and two BRDF material properties: Roughness and Metallicity. |
|
||||
| [prs-eth/marigold-iid-lighting-v1-1](https://huggingface.co/prs-eth/marigold-iid-lighting-v1-1) | Intrinsics | HyperSim decomposition of an image  \\(I\\)  is comprised of Albedo  \\(A\\), Diffuse shading  \\(S\\), and Non-diffuse residual  \\(R\\):  \\(I = A*S+R\\). |
|
||||
|
||||
<Tip>
|
||||
|
||||
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff
|
||||
between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to
|
||||
efficiently load the same components into multiple pipelines.
|
||||
Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section
|
||||
[here](../../using-diffusers/svd#reduce-memory-usage).
|
||||
|
||||
</Tip>
|
||||
|
||||
<Tip warning={true}>
|
||||
|
||||
Marigold pipelines were designed and tested only with `DDIMScheduler` and `LCMScheduler`.
|
||||
Depending on the scheduler, the number of inference steps required to get reliable predictions varies, and there is no universal value that works best across schedulers.
|
||||
Because of that, the default value of `num_inference_steps` in the `__call__` method of the pipeline is set to `None` (see the API reference).
|
||||
Unless set explicitly, its value will be taken from the checkpoint configuration `model_index.json`.
|
||||
This is done to ensure high-quality predictions when calling the pipeline with just the `image` argument.
|
||||
Marigold pipelines were designed and tested with the scheduler embedded in the model checkpoint.
|
||||
The optimal number of inference steps varies by scheduler, with no universal value that works best across all cases.
|
||||
To accommodate this, the `num_inference_steps` parameter in the pipeline's `__call__` method defaults to `None` (see the
|
||||
API reference).
|
||||
Unless set explicitly, it inherits the value from the `default_denoising_steps` field in the checkpoint configuration
|
||||
file (`model_index.json`).
|
||||
This ensures high-quality predictions when invoking the pipeline with only the `image` argument.
|
||||
|
||||
</Tip>
|
||||
|
||||
See also Marigold [usage examples](marigold_usage).
|
||||
See also Marigold [usage examples](../../using-diffusers/marigold_usage).
|
||||
|
||||
## Marigold Depth Prediction API
|
||||
|
||||
## MarigoldDepthPipeline
|
||||
[[autodoc]] MarigoldDepthPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## MarigoldNormalsPipeline
|
||||
[[autodoc]] MarigoldNormalsPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## MarigoldDepthOutput
|
||||
[[autodoc]] pipelines.marigold.pipeline_marigold_depth.MarigoldDepthOutput
|
||||
|
||||
## MarigoldNormalsOutput
|
||||
[[autodoc]] pipelines.marigold.pipeline_marigold_normals.MarigoldNormalsOutput
|
||||
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_depth
|
||||
|
||||
## Marigold Normals Estimation API
|
||||
[[autodoc]] MarigoldNormalsPipeline
|
||||
- __call__
|
||||
|
||||
[[autodoc]] pipelines.marigold.pipeline_marigold_normals.MarigoldNormalsOutput
|
||||
|
||||
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_normals
|
||||
|
||||
## Marigold Intrinsic Image Decomposition API
|
||||
|
||||
[[autodoc]] MarigoldIntrinsicsPipeline
|
||||
- __call__
|
||||
|
||||
[[autodoc]] pipelines.marigold.pipeline_marigold_intrinsics.MarigoldIntrinsicsOutput
|
||||
|
||||
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_intrinsics
|
||||
|
||||
@@ -15,6 +15,10 @@
|
||||
|
||||
# Mochi 1 Preview
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
> [!TIP]
|
||||
> Only a research preview of the model weights is available at the moment.
|
||||
|
||||
|
||||
80
docs/source/en/api/pipelines/omnigen.md
Normal file
80
docs/source/en/api/pipelines/omnigen.md
Normal file
@@ -0,0 +1,80 @@
|
||||
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License.
|
||||
-->
|
||||
|
||||
# OmniGen
|
||||
|
||||
[OmniGen: Unified Image Generation](https://arxiv.org/pdf/2409.11340) from BAAI, by Shitao Xiao, Yueze Wang, Junjie Zhou, Huaying Yuan, Xingrun Xing, Ruiran Yan, Chaofan Li, Shuting Wang, Tiejun Huang, Zheng Liu.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
*The emergence of Large Language Models (LLMs) has unified language generation tasks and revolutionized human-machine interaction. However, in the realm of image generation, a unified model capable of handling various tasks within a single framework remains largely unexplored. In this work, we introduce OmniGen, a new diffusion model for unified image generation. OmniGen is characterized by the following features: 1) Unification: OmniGen not only demonstrates text-to-image generation capabilities but also inherently supports various downstream tasks, such as image editing, subject-driven generation, and visual conditional generation. 2) Simplicity: The architecture of OmniGen is highly simplified, eliminating the need for additional plugins. Moreover, compared to existing diffusion models, it is more user-friendly and can complete complex tasks end-to-end through instructions without the need for extra intermediate steps, greatly simplifying the image generation workflow. 3) Knowledge Transfer: Benefit from learning in a unified format, OmniGen effectively transfers knowledge across different tasks, manages unseen tasks and domains, and exhibits novel capabilities. We also explore the model’s reasoning capabilities and potential applications of the chain-of-thought mechanism. This work represents the first attempt at a general-purpose image generation model, and we will release our resources at https://github.com/VectorSpaceLab/OmniGen to foster future advancements.*
|
||||
|
||||
<Tip>
|
||||
|
||||
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
|
||||
|
||||
</Tip>
|
||||
|
||||
This pipeline was contributed by [staoxiao](https://github.com/staoxiao). The original codebase can be found [here](https://github.com/VectorSpaceLab/OmniGen). The original weights can be found under [hf.co/shitao](https://huggingface.co/Shitao/OmniGen-v1).
|
||||
|
||||
## Inference
|
||||
|
||||
First, load the pipeline:
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import OmniGenPipeline
|
||||
|
||||
pipe = OmniGenPipeline.from_pretrained("Shitao/OmniGen-v1-diffusers", torch_dtype=torch.bfloat16)
|
||||
pipe.to("cuda")
|
||||
```
|
||||
|
||||
For text-to-image, pass a text prompt. By default, OmniGen generates a 1024x1024 image.
|
||||
You can try setting the `height` and `width` parameters to generate images with different size.
|
||||
|
||||
```python
|
||||
prompt = "Realistic photo. A young woman sits on a sofa, holding a book and facing the camera. She wears delicate silver hoop earrings adorned with tiny, sparkling diamonds that catch the light, with her long chestnut hair cascading over her shoulders. Her eyes are focused and gentle, framed by long, dark lashes. She is dressed in a cozy cream sweater, which complements her warm, inviting smile. Behind her, there is a table with a cup of water in a sleek, minimalist blue mug. The background is a serene indoor setting with soft natural light filtering through a window, adorned with tasteful art and flowers, creating a cozy and peaceful ambiance. 4K, HD."
|
||||
image = pipe(
|
||||
prompt=prompt,
|
||||
height=1024,
|
||||
width=1024,
|
||||
guidance_scale=3,
|
||||
generator=torch.Generator(device="cpu").manual_seed(111),
|
||||
).images[0]
|
||||
image.save("output.png")
|
||||
```
|
||||
|
||||
OmniGen supports multimodal inputs.
|
||||
When the input includes an image, you need to add a placeholder `<img><|image_1|></img>` in the text prompt to represent the image.
|
||||
It is recommended to enable `use_input_image_size_as_output` to keep the edited image the same size as the original image.
|
||||
|
||||
```python
|
||||
prompt="<img><|image_1|></img> Remove the woman's earrings. Replace the mug with a clear glass filled with sparkling iced cola."
|
||||
input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png")]
|
||||
image = pipe(
|
||||
prompt=prompt,
|
||||
input_images=input_images,
|
||||
guidance_scale=2,
|
||||
img_guidance_scale=1.6,
|
||||
use_input_image_size_as_output=True,
|
||||
generator=torch.Generator(device="cpu").manual_seed(222)).images[0]
|
||||
image.save("output.png")
|
||||
```
|
||||
|
||||
## OmniGenPipeline
|
||||
|
||||
[[autodoc]] OmniGenPipeline
|
||||
- all
|
||||
- __call__
|
||||
@@ -54,7 +54,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
|
||||
| [DiT](dit) | text2image |
|
||||
| [Flux](flux) | text2image |
|
||||
| [Hunyuan-DiT](hunyuandit) | text2image |
|
||||
| [I2VGen-XL](i2vgenxl) | text2video |
|
||||
| [I2VGen-XL](i2vgenxl) | image2video |
|
||||
| [InstructPix2Pix](pix2pix) | image editing |
|
||||
| [Kandinsky 2.1](kandinsky) | text2image, image2image, inpainting, interpolation |
|
||||
| [Kandinsky 2.2](kandinsky_v22) | text2image, image2image, inpainting |
|
||||
@@ -65,7 +65,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
|
||||
| [Latte](latte) | text2image |
|
||||
| [LEDITS++](ledits_pp) | image editing |
|
||||
| [Lumina-T2X](lumina) | text2image |
|
||||
| [Marigold](marigold) | depth |
|
||||
| [Marigold](marigold) | depth-estimation, normals-estimation, intrinsic-decomposition |
|
||||
| [MultiDiffusion](panorama) | text2image |
|
||||
| [MusicLDM](musicldm) | text2audio |
|
||||
| [PAG](pag) | text2image |
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Perturbed-Attention Guidance
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[Perturbed-Attention Guidance (PAG)](https://ku-cvlab.github.io/Perturbed-Attention-Guidance/) is a new diffusion sampling guidance that improves sample quality across both unconditional and conditional settings, achieving this without requiring further training or the integration of external modules.
|
||||
|
||||
PAG was introduced in [Self-Rectifying Diffusion Sampling with Perturbed-Attention Guidance](https://huggingface.co/papers/2403.17377) by Donghoon Ahn, Hyoungwon Cho, Jaewon Min, Wooseok Jang, Jungwoo Kim, SeonHwa Kim, Hyun Hee Park, Kyong Hwan Jin and Seungryong Kim.
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# MultiDiffusion
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation](https://huggingface.co/papers/2302.08113) is by Omer Bar-Tal, Lior Yariv, Yaron Lipman, and Tali Dekel.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Image-to-Video Generation with PIA (Personalized Image Animator)
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
## Overview
|
||||
|
||||
[PIA: Your Personalized Image Animator via Plug-and-Play Modules in Text-to-Image Models](https://arxiv.org/abs/2312.13964) by Yiming Zhang, Zhening Xing, Yanhong Zeng, Youqing Fang, Kai Chen
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# InstructPix2Pix
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[InstructPix2Pix: Learning to Follow Image Editing Instructions](https://huggingface.co/papers/2211.09800) is by Tim Brooks, Aleksander Holynski and Alexei A. Efros.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
@@ -14,6 +14,11 @@
|
||||
|
||||
# SanaPipeline
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
|
||||
</div>
|
||||
|
||||
[SANA: Efficient High-Resolution Image Synthesis with Linear Diffusion Transformers](https://huggingface.co/papers/2410.10629) from NVIDIA and MIT HAN Lab, by Enze Xie, Junsong Chen, Junyu Chen, Han Cai, Haotian Tang, Yujun Lin, Zhekai Zhang, Muyang Li, Ligeng Zhu, Yao Lu, Song Han.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
100
docs/source/en/api/pipelines/sana_sprint.md
Normal file
100
docs/source/en/api/pipelines/sana_sprint.md
Normal file
@@ -0,0 +1,100 @@
|
||||
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License. -->
|
||||
|
||||
# SANA-Sprint
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[SANA-Sprint: One-Step Diffusion with Continuous-Time Consistency Distillation](https://huggingface.co/papers/2503.09641) from NVIDIA, MIT HAN Lab, and Hugging Face by Junsong Chen, Shuchen Xue, Yuyang Zhao, Jincheng Yu, Sayak Paul, Junyu Chen, Han Cai, Enze Xie, Song Han
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
*This paper presents SANA-Sprint, an efficient diffusion model for ultra-fast text-to-image (T2I) generation. SANA-Sprint is built on a pre-trained foundation model and augmented with hybrid distillation, dramatically reducing inference steps from 20 to 1-4. We introduce three key innovations: (1) We propose a training-free approach that transforms a pre-trained flow-matching model for continuous-time consistency distillation (sCM), eliminating costly training from scratch and achieving high training efficiency. Our hybrid distillation strategy combines sCM with latent adversarial distillation (LADD): sCM ensures alignment with the teacher model, while LADD enhances single-step generation fidelity. (2) SANA-Sprint is a unified step-adaptive model that achieves high-quality generation in 1-4 steps, eliminating step-specific training and improving efficiency. (3) We integrate ControlNet with SANA-Sprint for real-time interactive image generation, enabling instant visual feedback for user interaction. SANA-Sprint establishes a new Pareto frontier in speed-quality tradeoffs, achieving state-of-the-art performance with 7.59 FID and 0.74 GenEval in only 1 step — outperforming FLUX-schnell (7.94 FID / 0.71 GenEval) while being 10× faster (0.1s vs 1.1s on H100). It also achieves 0.1s (T2I) and 0.25s (ControlNet) latency for 1024×1024 images on H100, and 0.31s (T2I) on an RTX 4090, showcasing its exceptional efficiency and potential for AI-powered consumer applications (AIPC). Code and pre-trained models will be open-sourced.*
|
||||
|
||||
<Tip>
|
||||
|
||||
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
|
||||
|
||||
</Tip>
|
||||
|
||||
This pipeline was contributed by [lawrence-cj](https://github.com/lawrence-cj), [shuchen Xue](https://github.com/scxue) and [Enze Xie](https://github.com/xieenze). The original codebase can be found [here](https://github.com/NVlabs/Sana). The original weights can be found under [hf.co/Efficient-Large-Model](https://huggingface.co/Efficient-Large-Model/).
|
||||
|
||||
Available models:
|
||||
|
||||
| Model | Recommended dtype |
|
||||
|:-------------------------------------------------------------------------------------------------------------------------------------------:|:-----------------:|
|
||||
| [`Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers) | `torch.bfloat16` |
|
||||
| [`Efficient-Large-Model/Sana_Sprint_0.6B_1024px_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_Sprint_0.6B_1024px_diffusers) | `torch.bfloat16` |
|
||||
|
||||
Refer to [this](https://huggingface.co/collections/Efficient-Large-Model/sana-sprint-67d6810d65235085b3b17c76) collection for more information.
|
||||
|
||||
Note: The recommended dtype mentioned is for the transformer weights. The text encoder must stay in `torch.bfloat16` and VAE weights must stay in `torch.bfloat16` or `torch.float32` for the model to work correctly. Please refer to the inference example below to see how to load the model with the recommended dtype.
|
||||
|
||||
|
||||
## Quantization
|
||||
|
||||
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
|
||||
|
||||
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`SanaSprintPipeline`] for inference with bitsandbytes.
|
||||
|
||||
```py
|
||||
import torch
|
||||
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, SanaTransformer2DModel, SanaSprintPipeline
|
||||
from transformers import BitsAndBytesConfig as BitsAndBytesConfig, AutoModel
|
||||
|
||||
quant_config = BitsAndBytesConfig(load_in_8bit=True)
|
||||
text_encoder_8bit = AutoModel.from_pretrained(
|
||||
"Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers",
|
||||
subfolder="text_encoder",
|
||||
quantization_config=quant_config,
|
||||
torch_dtype=torch.bfloat16,
|
||||
)
|
||||
|
||||
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True)
|
||||
transformer_8bit = SanaTransformer2DModel.from_pretrained(
|
||||
"Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers",
|
||||
subfolder="transformer",
|
||||
quantization_config=quant_config,
|
||||
torch_dtype=torch.bfloat16,
|
||||
)
|
||||
|
||||
pipeline = SanaSprintPipeline.from_pretrained(
|
||||
"Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers",
|
||||
text_encoder=text_encoder_8bit,
|
||||
transformer=transformer_8bit,
|
||||
torch_dtype=torch.bfloat16,
|
||||
device_map="balanced",
|
||||
)
|
||||
|
||||
prompt = "a tiny astronaut hatching from an egg on the moon"
|
||||
image = pipeline(prompt).images[0]
|
||||
image.save("sana.png")
|
||||
```
|
||||
|
||||
## Setting `max_timesteps`
|
||||
|
||||
Users can tweak the `max_timesteps` value for experimenting with the visual quality of the generated outputs. The default `max_timesteps` value was obtained with an inference-time search process. For more details about it, check out the paper.
|
||||
|
||||
## SanaSprintPipeline
|
||||
|
||||
[[autodoc]] SanaSprintPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
|
||||
## SanaPipelineOutput
|
||||
|
||||
[[autodoc]] pipelines.sana.pipeline_output.SanaPipelineOutput
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Depth-to-image
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
The Stable Diffusion model can also infer depth based on an image using [MiDaS](https://github.com/isl-org/MiDaS). This allows you to pass a text prompt and an initial image to condition the generation of new images as well as a `depth_map` to preserve the image structure.
|
||||
|
||||
<Tip>
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Image-to-image
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
The Stable Diffusion model can also be applied to image-to-image generation by passing a text prompt and an initial image to condition the generation of new images.
|
||||
|
||||
The [`StableDiffusionImg2ImgPipeline`] uses the diffusion-denoising mechanism proposed in [SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equations](https://huggingface.co/papers/2108.01073) by Chenlin Meng, Yutong He, Yang Song, Jiaming Song, Jiajun Wu, Jun-Yan Zhu, Stefano Ermon.
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Inpainting
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
The Stable Diffusion model can also be applied to inpainting which lets you edit specific parts of an image by providing a mask and a text prompt using Stable Diffusion.
|
||||
|
||||
## Tips
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Text-to-(RGB, depth)
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
LDM3D was proposed in [LDM3D: Latent Diffusion Model for 3D](https://huggingface.co/papers/2305.10853) by Gabriela Ben Melech Stan, Diana Wofk, Scottie Fox, Alex Redden, Will Saxton, Jean Yu, Estelle Aflalo, Shao-Yen Tseng, Fabio Nonato, Matthias Muller, and Vasudev Lal. LDM3D generates an image and a depth map from a given text prompt unlike the existing text-to-image diffusion models such as [Stable Diffusion](./overview) which only generates an image. With almost the same number of parameters, LDM3D achieves to create a latent space that can compress both the RGB images and the depth maps.
|
||||
|
||||
Two checkpoints are available for use:
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Stable Diffusion pipelines
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/). Latent diffusion applies the diffusion process over a lower dimensional latent space to reduce memory and compute complexity. This specific type of diffusion model was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
|
||||
|
||||
Stable Diffusion is trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs.
|
||||
|
||||
@@ -12,6 +12,11 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Stable Diffusion 3
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
|
||||
</div>
|
||||
|
||||
Stable Diffusion 3 (SD3) was proposed in [Scaling Rectified Flow Transformers for High-Resolution Image Synthesis](https://arxiv.org/pdf/2403.03206.pdf) by Patrick Esser, Sumith Kulal, Andreas Blattmann, Rahim Entezari, Jonas Muller, Harry Saini, Yam Levi, Dominik Lorenz, Axel Sauer, Frederic Boesel, Dustin Podell, Tim Dockhorn, Zion English, Kyle Lacey, Alex Goodwin, Yannik Marek, and Robin Rombach.
|
||||
|
||||
The abstract from the paper is:
|
||||
@@ -77,7 +82,7 @@ from diffusers import StableDiffusion3Pipeline
|
||||
from transformers import SiglipVisionModel, SiglipImageProcessor
|
||||
|
||||
image_encoder_id = "google/siglip-so400m-patch14-384"
|
||||
ip_adapter_id = "guiyrt/InstantX-SD3.5-Large-IP-Adapter-diffusers"
|
||||
ip_adapter_id = "InstantX/SD3.5-Large-IP-Adapter"
|
||||
|
||||
feature_extractor = SiglipImageProcessor.from_pretrained(
|
||||
image_encoder_id,
|
||||
|
||||
@@ -12,6 +12,11 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Stable Diffusion XL
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
|
||||
</div>
|
||||
|
||||
Stable Diffusion XL (SDXL) was proposed in [SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis](https://huggingface.co/papers/2307.01952) by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, and Robin Rombach.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Text-to-image
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
The Stable Diffusion model was created by researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [Runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionPipeline`] is capable of generating photorealistic images given any text input. It's trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs. Latent diffusion is the research on top of which Stable Diffusion was built. It was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Super-resolution
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
The Stable Diffusion upscaler diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), and [LAION](https://laion.ai/). It is used to enhance the resolution of input images by a factor of 4.
|
||||
|
||||
<Tip>
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Stable unCLIP
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
Stable unCLIP checkpoints are finetuned from [Stable Diffusion 2.1](./stable_diffusion/stable_diffusion_2) checkpoints to condition on CLIP image embeddings.
|
||||
Stable unCLIP still conditions on text embeddings. Given the two separate conditionings, stable unCLIP can be used
|
||||
for text guided image variation. When combined with an unCLIP prior, it can also be used for full text to image generation.
|
||||
|
||||
@@ -18,6 +18,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Text-to-video
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[ModelScope Text-to-Video Technical Report](https://arxiv.org/abs/2308.06571) is by Jiuniu Wang, Hangjie Yuan, Dayou Chen, Yingya Zhang, Xiang Wang, Shiwei Zhang.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Text2Video-Zero
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[Text2Video-Zero: Text-to-Image Diffusion Models are Zero-Shot Video Generators](https://huggingface.co/papers/2303.13439) is by Levon Khachatryan, Andranik Movsisyan, Vahram Tadevosyan, Roberto Henschel, [Zhangyang Wang](https://www.ece.utexas.edu/people/faculty/atlas-wang), Shant Navasardyan, [Humphrey Shi](https://www.humphreyshi.com).
|
||||
|
||||
Text2Video-Zero enables zero-shot video generation using either:
|
||||
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# UniDiffuser
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
The UniDiffuser model was proposed in [One Transformer Fits All Distributions in Multi-Modal Diffusion at Scale](https://huggingface.co/papers/2303.06555) by Fan Bao, Shen Nie, Kaiwen Xue, Chongxuan Li, Shi Pu, Yaole Wang, Gang Yue, Yue Cao, Hang Su, Jun Zhu.
|
||||
|
||||
The abstract from the paper is:
|
||||
|
||||
519
docs/source/en/api/pipelines/wan.md
Normal file
519
docs/source/en/api/pipelines/wan.md
Normal file
@@ -0,0 +1,519 @@
|
||||
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License. -->
|
||||
|
||||
# Wan
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
[Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
|
||||
|
||||
<!-- TODO(aryan): update abstract once paper is out -->
|
||||
|
||||
## Generating Videos with Wan 2.1
|
||||
|
||||
We will first need to install some addtional dependencies.
|
||||
|
||||
```shell
|
||||
pip install -u ftfy imageio-ffmpeg imageio
|
||||
```
|
||||
|
||||
### Text to Video Generation
|
||||
|
||||
The following example requires 11GB VRAM to run and uses the smaller `Wan-AI/Wan2.1-T2V-1.3B-Diffusers` model. You can switch it out
|
||||
for the larger `Wan2.1-I2V-14B-720P-Diffusers` or `Wan-AI/Wan2.1-I2V-14B-480P-Diffusers` if you have at least 35GB VRAM available.
|
||||
|
||||
```python
|
||||
from diffusers import WanPipeline
|
||||
from diffusers.utils import export_to_video
|
||||
|
||||
# Available models: Wan-AI/Wan2.1-I2V-14B-720P-Diffusers or Wan-AI/Wan2.1-I2V-14B-480P-Diffusers
|
||||
model_id = "Wan-AI/Wan2.1-T2V-1.3B-Diffusers"
|
||||
|
||||
pipe = WanPipeline.from_pretrained(model_id, torch_dtype=torch.bfloat16)
|
||||
pipe.enable_model_cpu_offload()
|
||||
|
||||
prompt = "A cat and a dog baking a cake together in a kitchen. The cat is carefully measuring flour, while the dog is stirring the batter with a wooden spoon. The kitchen is cozy, with sunlight streaming through the window."
|
||||
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
|
||||
num_frames = 33
|
||||
|
||||
frames = pipe(prompt=prompt, negative_prompt=negative_prompt, num_frames=num_frames).frames[0]
|
||||
export_to_video(frames, "wan-t2v.mp4", fps=16)
|
||||
```
|
||||
|
||||
<Tip>
|
||||
You can improve the quality of the generated video by running the decoding step in full precision.
|
||||
</Tip>
|
||||
|
||||
```python
|
||||
from diffusers import WanPipeline, AutoencoderKLWan
|
||||
from diffusers.utils import export_to_video
|
||||
|
||||
model_id = "Wan-AI/Wan2.1-T2V-1.3B-Diffusers"
|
||||
|
||||
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
|
||||
pipe = WanPipeline.from_pretrained(model_id, vae=vae, torch_dtype=torch.bfloat16)
|
||||
|
||||
# replace this with pipe.to("cuda") if you have sufficient VRAM
|
||||
pipe.enable_model_cpu_offload()
|
||||
|
||||
prompt = "A cat and a dog baking a cake together in a kitchen. The cat is carefully measuring flour, while the dog is stirring the batter with a wooden spoon. The kitchen is cozy, with sunlight streaming through the window."
|
||||
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
|
||||
num_frames = 33
|
||||
|
||||
frames = pipe(prompt=prompt, num_frames=num_frames).frames[0]
|
||||
export_to_video(frames, "wan-t2v.mp4", fps=16)
|
||||
```
|
||||
|
||||
### Image to Video Generation
|
||||
|
||||
The Image to Video pipeline requires loading the `AutoencoderKLWan` and the `CLIPVisionModel` components in full precision. The following example will need at least
|
||||
35GB of VRAM to run.
|
||||
|
||||
```python
|
||||
import torch
|
||||
import numpy as np
|
||||
from diffusers import AutoencoderKLWan, WanImageToVideoPipeline
|
||||
from diffusers.utils import export_to_video, load_image
|
||||
from transformers import CLIPVisionModel
|
||||
|
||||
# Available models: Wan-AI/Wan2.1-I2V-14B-480P-Diffusers, Wan-AI/Wan2.1-I2V-14B-720P-Diffusers
|
||||
model_id = "Wan-AI/Wan2.1-I2V-14B-480P-Diffusers"
|
||||
image_encoder = CLIPVisionModel.from_pretrained(
|
||||
model_id, subfolder="image_encoder", torch_dtype=torch.float32
|
||||
)
|
||||
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
|
||||
pipe = WanImageToVideoPipeline.from_pretrained(
|
||||
model_id, vae=vae, image_encoder=image_encoder, torch_dtype=torch.bfloat16
|
||||
)
|
||||
|
||||
# replace this with pipe.to("cuda") if you have sufficient VRAM
|
||||
pipe.enable_model_cpu_offload()
|
||||
|
||||
image = load_image(
|
||||
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"
|
||||
)
|
||||
|
||||
max_area = 480 * 832
|
||||
aspect_ratio = image.height / image.width
|
||||
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
|
||||
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
|
||||
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
|
||||
image = image.resize((width, height))
|
||||
|
||||
prompt = (
|
||||
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
|
||||
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
|
||||
)
|
||||
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
|
||||
|
||||
num_frames = 33
|
||||
|
||||
output = pipe(
|
||||
image=image,
|
||||
prompt=prompt,
|
||||
negative_prompt=negative_prompt,
|
||||
height=height,
|
||||
width=width,
|
||||
num_frames=num_frames,
|
||||
guidance_scale=5.0,
|
||||
).frames[0]
|
||||
export_to_video(output, "wan-i2v.mp4", fps=16)
|
||||
```
|
||||
|
||||
### First and Last Frame Interpolation
|
||||
|
||||
```python
|
||||
import numpy as np
|
||||
import torch
|
||||
import torchvision.transforms.functional as TF
|
||||
from diffusers import AutoencoderKLWan, WanImageToVideoPipeline
|
||||
from diffusers.utils import export_to_video, load_image
|
||||
from transformers import CLIPVisionModel
|
||||
|
||||
|
||||
model_id = "Wan-AI/Wan2.1-FLF2V-14B-720P-diffusers"
|
||||
image_encoder = CLIPVisionModel.from_pretrained(model_id, subfolder="image_encoder", torch_dtype=torch.float32)
|
||||
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
|
||||
pipe = WanImageToVideoPipeline.from_pretrained(
|
||||
model_id, vae=vae, image_encoder=image_encoder, torch_dtype=torch.bfloat16
|
||||
)
|
||||
pipe.to("cuda")
|
||||
|
||||
first_frame = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flf2v_input_first_frame.png")
|
||||
last_frame = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flf2v_input_last_frame.png")
|
||||
|
||||
def aspect_ratio_resize(image, pipe, max_area=720 * 1280):
|
||||
aspect_ratio = image.height / image.width
|
||||
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
|
||||
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
|
||||
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
|
||||
image = image.resize((width, height))
|
||||
return image, height, width
|
||||
|
||||
def center_crop_resize(image, height, width):
|
||||
# Calculate resize ratio to match first frame dimensions
|
||||
resize_ratio = max(width / image.width, height / image.height)
|
||||
|
||||
# Resize the image
|
||||
width = round(image.width * resize_ratio)
|
||||
height = round(image.height * resize_ratio)
|
||||
size = [width, height]
|
||||
image = TF.center_crop(image, size)
|
||||
|
||||
return image, height, width
|
||||
|
||||
first_frame, height, width = aspect_ratio_resize(first_frame, pipe)
|
||||
if last_frame.size != first_frame.size:
|
||||
last_frame, _, _ = center_crop_resize(last_frame, height, width)
|
||||
|
||||
prompt = "CG animation style, a small blue bird takes off from the ground, flapping its wings. The bird's feathers are delicate, with a unique pattern on its chest. The background shows a blue sky with white clouds under bright sunshine. The camera follows the bird upward, capturing its flight and the vastness of the sky from a close-up, low-angle perspective."
|
||||
|
||||
output = pipe(
|
||||
image=first_frame, last_image=last_frame, prompt=prompt, height=height, width=width, guidance_scale=5.5
|
||||
).frames[0]
|
||||
export_to_video(output, "output.mp4", fps=16)
|
||||
```
|
||||
|
||||
### Video to Video Generation
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers.utils import load_video, export_to_video
|
||||
from diffusers import AutoencoderKLWan, WanVideoToVideoPipeline, UniPCMultistepScheduler
|
||||
|
||||
# Available models: Wan-AI/Wan2.1-T2V-14B-Diffusers, Wan-AI/Wan2.1-T2V-1.3B-Diffusers
|
||||
model_id = "Wan-AI/Wan2.1-T2V-1.3B-Diffusers"
|
||||
vae = AutoencoderKLWan.from_pretrained(
|
||||
model_id, subfolder="vae", torch_dtype=torch.float32
|
||||
)
|
||||
pipe = WanVideoToVideoPipeline.from_pretrained(
|
||||
model_id, vae=vae, torch_dtype=torch.bfloat16
|
||||
)
|
||||
flow_shift = 3.0 # 5.0 for 720P, 3.0 for 480P
|
||||
pipe.scheduler = UniPCMultistepScheduler.from_config(
|
||||
pipe.scheduler.config, flow_shift=flow_shift
|
||||
)
|
||||
# change to pipe.to("cuda") if you have sufficient VRAM
|
||||
pipe.enable_model_cpu_offload()
|
||||
|
||||
prompt = "A robot standing on a mountain top. The sun is setting in the background"
|
||||
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
|
||||
video = load_video(
|
||||
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/hiker.mp4"
|
||||
)
|
||||
output = pipe(
|
||||
video=video,
|
||||
prompt=prompt,
|
||||
negative_prompt=negative_prompt,
|
||||
height=480,
|
||||
width=512,
|
||||
guidance_scale=7.0,
|
||||
strength=0.7,
|
||||
).frames[0]
|
||||
|
||||
export_to_video(output, "wan-v2v.mp4", fps=16)
|
||||
```
|
||||
|
||||
## Memory Optimizations for Wan 2.1
|
||||
|
||||
Base inference with the large 14B Wan 2.1 models can take up to 35GB of VRAM when generating videos at 720p resolution. We'll outline a few memory optimizations we can apply to reduce the VRAM required to run the model.
|
||||
|
||||
We'll use `Wan-AI/Wan2.1-I2V-14B-720P-Diffusers` model in these examples to demonstrate the memory savings, but the techniques are applicable to all model checkpoints.
|
||||
|
||||
### Group Offloading the Transformer and UMT5 Text Encoder
|
||||
|
||||
Find more information about group offloading [here](../optimization/memory.md)
|
||||
|
||||
#### Block Level Group Offloading
|
||||
|
||||
We can reduce our VRAM requirements by applying group offloading to the larger model components of the pipeline; the `WanTransformer3DModel` and `UMT5EncoderModel`. Group offloading will break up the individual modules of a model and offload/onload them onto your GPU as needed during inference. In this example, we'll apply `block_level` offloading, which will group the modules in a model into blocks of size `num_blocks_per_group` and offload/onload them to GPU. Moving to between CPU and GPU does add latency to the inference process. You can trade off between latency and memory savings by increasing or decreasing the `num_blocks_per_group`.
|
||||
|
||||
The following example will now only require 14GB of VRAM to run, but will take approximately 30 minutes to generate a video.
|
||||
|
||||
```python
|
||||
import torch
|
||||
import numpy as np
|
||||
from diffusers import AutoencoderKLWan, WanTransformer3DModel, WanImageToVideoPipeline
|
||||
from diffusers.hooks.group_offloading import apply_group_offloading
|
||||
from diffusers.utils import export_to_video, load_image
|
||||
from transformers import UMT5EncoderModel, CLIPVisionModel
|
||||
|
||||
# Available models: Wan-AI/Wan2.1-I2V-14B-480P-Diffusers, Wan-AI/Wan2.1-I2V-14B-720P-Diffusers
|
||||
model_id = "Wan-AI/Wan2.1-I2V-14B-720P-Diffusers"
|
||||
image_encoder = CLIPVisionModel.from_pretrained(
|
||||
model_id, subfolder="image_encoder", torch_dtype=torch.float32
|
||||
)
|
||||
|
||||
text_encoder = UMT5EncoderModel.from_pretrained(model_id, subfolder="text_encoder", torch_dtype=torch.bfloat16)
|
||||
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
|
||||
transformer = WanTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
|
||||
|
||||
onload_device = torch.device("cuda")
|
||||
offload_device = torch.device("cpu")
|
||||
|
||||
apply_group_offloading(text_encoder,
|
||||
onload_device=onload_device,
|
||||
offload_device=offload_device,
|
||||
offload_type="block_level",
|
||||
num_blocks_per_group=4
|
||||
)
|
||||
|
||||
transformer.enable_group_offload(
|
||||
onload_device=onload_device,
|
||||
offload_device=offload_device,
|
||||
offload_type="block_level",
|
||||
num_blocks_per_group=4,
|
||||
)
|
||||
pipe = WanImageToVideoPipeline.from_pretrained(
|
||||
model_id,
|
||||
vae=vae,
|
||||
transformer=transformer,
|
||||
text_encoder=text_encoder,
|
||||
image_encoder=image_encoder,
|
||||
torch_dtype=torch.bfloat16
|
||||
)
|
||||
# Since we've offloaded the larger models alrady, we can move the rest of the model components to GPU
|
||||
pipe.to("cuda")
|
||||
|
||||
image = load_image(
|
||||
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"
|
||||
)
|
||||
|
||||
max_area = 720 * 832
|
||||
aspect_ratio = image.height / image.width
|
||||
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
|
||||
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
|
||||
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
|
||||
image = image.resize((width, height))
|
||||
|
||||
prompt = (
|
||||
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
|
||||
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
|
||||
)
|
||||
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
|
||||
|
||||
num_frames = 33
|
||||
|
||||
output = pipe(
|
||||
image=image,
|
||||
prompt=prompt,
|
||||
negative_prompt=negative_prompt,
|
||||
height=height,
|
||||
width=width,
|
||||
num_frames=num_frames,
|
||||
guidance_scale=5.0,
|
||||
).frames[0]
|
||||
|
||||
export_to_video(output, "wan-i2v.mp4", fps=16)
|
||||
```
|
||||
|
||||
#### Block Level Group Offloading with CUDA Streams
|
||||
|
||||
We can speed up group offloading inference, by enabling the use of [CUDA streams](https://pytorch.org/docs/stable/generated/torch.cuda.Stream.html). However, using CUDA streams requires moving the model parameters into pinned memory. This allocation is handled by Pytorch under the hood, and can result in a significant spike in CPU RAM usage. Please consider this option if your CPU RAM is atleast 2X the size of the model you are group offloading.
|
||||
|
||||
In the following example we will use CUDA streams when group offloading the `WanTransformer3DModel`. When testing on an A100, this example will require 14GB of VRAM, 52GB of CPU RAM, but will generate a video in approximately 9 minutes.
|
||||
|
||||
```python
|
||||
import torch
|
||||
import numpy as np
|
||||
from diffusers import AutoencoderKLWan, WanTransformer3DModel, WanImageToVideoPipeline
|
||||
from diffusers.hooks.group_offloading import apply_group_offloading
|
||||
from diffusers.utils import export_to_video, load_image
|
||||
from transformers import UMT5EncoderModel, CLIPVisionModel
|
||||
|
||||
# Available models: Wan-AI/Wan2.1-I2V-14B-480P-Diffusers, Wan-AI/Wan2.1-I2V-14B-720P-Diffusers
|
||||
model_id = "Wan-AI/Wan2.1-I2V-14B-720P-Diffusers"
|
||||
image_encoder = CLIPVisionModel.from_pretrained(
|
||||
model_id, subfolder="image_encoder", torch_dtype=torch.float32
|
||||
)
|
||||
|
||||
text_encoder = UMT5EncoderModel.from_pretrained(model_id, subfolder="text_encoder", torch_dtype=torch.bfloat16)
|
||||
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
|
||||
transformer = WanTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
|
||||
|
||||
onload_device = torch.device("cuda")
|
||||
offload_device = torch.device("cpu")
|
||||
|
||||
apply_group_offloading(text_encoder,
|
||||
onload_device=onload_device,
|
||||
offload_device=offload_device,
|
||||
offload_type="block_level",
|
||||
num_blocks_per_group=4
|
||||
)
|
||||
|
||||
transformer.enable_group_offload(
|
||||
onload_device=onload_device,
|
||||
offload_device=offload_device,
|
||||
offload_type="leaf_level",
|
||||
use_stream=True
|
||||
)
|
||||
pipe = WanImageToVideoPipeline.from_pretrained(
|
||||
model_id,
|
||||
vae=vae,
|
||||
transformer=transformer,
|
||||
text_encoder=text_encoder,
|
||||
image_encoder=image_encoder,
|
||||
torch_dtype=torch.bfloat16
|
||||
)
|
||||
# Since we've offloaded the larger models alrady, we can move the rest of the model components to GPU
|
||||
pipe.to("cuda")
|
||||
|
||||
image = load_image(
|
||||
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"
|
||||
)
|
||||
|
||||
max_area = 720 * 832
|
||||
aspect_ratio = image.height / image.width
|
||||
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
|
||||
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
|
||||
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
|
||||
image = image.resize((width, height))
|
||||
|
||||
prompt = (
|
||||
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
|
||||
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
|
||||
)
|
||||
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
|
||||
|
||||
num_frames = 33
|
||||
|
||||
output = pipe(
|
||||
image=image,
|
||||
prompt=prompt,
|
||||
negative_prompt=negative_prompt,
|
||||
height=height,
|
||||
width=width,
|
||||
num_frames=num_frames,
|
||||
guidance_scale=5.0,
|
||||
).frames[0]
|
||||
|
||||
export_to_video(output, "wan-i2v.mp4", fps=16)
|
||||
```
|
||||
|
||||
### Applying Layerwise Casting to the Transformer
|
||||
|
||||
Find more information about layerwise casting [here](../optimization/memory.md)
|
||||
|
||||
In this example, we will model offloading with layerwise casting. Layerwise casting will downcast each layer's weights to `torch.float8_e4m3fn`, temporarily upcast to `torch.bfloat16` during the forward pass of the layer, then revert to `torch.float8_e4m3fn` afterward. This approach reduces memory requirements by approximately 50% while introducing a minor quality reduction in the generated video due to the precision trade-off.
|
||||
|
||||
This example will require 20GB of VRAM.
|
||||
|
||||
```python
|
||||
import torch
|
||||
import numpy as np
|
||||
from diffusers import AutoencoderKLWan, WanTransformer3DModel, WanImageToVideoPipeline
|
||||
from diffusers.hooks.group_offloading import apply_group_offloading
|
||||
from diffusers.utils import export_to_video, load_image
|
||||
from transformers import UMT5EncoderModel, CLIPVisionModel
|
||||
|
||||
model_id = "Wan-AI/Wan2.1-I2V-14B-720P-Diffusers"
|
||||
image_encoder = CLIPVisionModel.from_pretrained(
|
||||
model_id, subfolder="image_encoder", torch_dtype=torch.float32
|
||||
)
|
||||
text_encoder = UMT5EncoderModel.from_pretrained(model_id, subfolder="text_encoder", torch_dtype=torch.bfloat16)
|
||||
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
|
||||
|
||||
transformer = WanTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
|
||||
transformer.enable_layerwise_casting(storage_dtype=torch.float8_e4m3fn, compute_dtype=torch.bfloat16)
|
||||
|
||||
pipe = WanImageToVideoPipeline.from_pretrained(
|
||||
model_id,
|
||||
vae=vae,
|
||||
transformer=transformer,
|
||||
text_encoder=text_encoder,
|
||||
image_encoder=image_encoder,
|
||||
torch_dtype=torch.bfloat16
|
||||
)
|
||||
pipe.enable_model_cpu_offload()
|
||||
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg")
|
||||
|
||||
max_area = 720 * 832
|
||||
aspect_ratio = image.height / image.width
|
||||
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
|
||||
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
|
||||
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
|
||||
image = image.resize((width, height))
|
||||
prompt = (
|
||||
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
|
||||
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
|
||||
)
|
||||
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
|
||||
num_frames = 33
|
||||
|
||||
output = pipe(
|
||||
image=image,
|
||||
prompt=prompt,
|
||||
negative_prompt=negative_prompt,
|
||||
height=height,
|
||||
width=width,
|
||||
num_frames=num_frames,
|
||||
num_inference_steps=50,
|
||||
guidance_scale=5.0,
|
||||
).frames[0]
|
||||
export_to_video(output, "wan-i2v.mp4", fps=16)
|
||||
```
|
||||
|
||||
## Using a Custom Scheduler
|
||||
|
||||
Wan can be used with many different schedulers, each with their own benefits regarding speed and generation quality. By default, Wan uses the `UniPCMultistepScheduler(prediction_type="flow_prediction", use_flow_sigmas=True, flow_shift=3.0)` scheduler. You can use a different scheduler as follows:
|
||||
|
||||
```python
|
||||
from diffusers import FlowMatchEulerDiscreteScheduler, UniPCMultistepScheduler, WanPipeline
|
||||
|
||||
scheduler_a = FlowMatchEulerDiscreteScheduler(shift=5.0)
|
||||
scheduler_b = UniPCMultistepScheduler(prediction_type="flow_prediction", use_flow_sigmas=True, flow_shift=4.0)
|
||||
|
||||
pipe = WanPipeline.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", scheduler=<CUSTOM_SCHEDULER_HERE>)
|
||||
|
||||
# or,
|
||||
pipe.scheduler = <CUSTOM_SCHEDULER_HERE>
|
||||
```
|
||||
|
||||
## Using Single File Loading with Wan 2.1
|
||||
|
||||
The `WanTransformer3DModel` and `AutoencoderKLWan` models support loading checkpoints in their original format via the `from_single_file` loading
|
||||
method.
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import WanPipeline, WanTransformer3DModel
|
||||
|
||||
ckpt_path = "https://huggingface.co/Comfy-Org/Wan_2.1_ComfyUI_repackaged/blob/main/split_files/diffusion_models/wan2.1_t2v_1.3B_bf16.safetensors"
|
||||
transformer = WanTransformer3DModel.from_single_file(ckpt_path, torch_dtype=torch.bfloat16)
|
||||
|
||||
pipe = WanPipeline.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", transformer=transformer)
|
||||
```
|
||||
|
||||
## Recommendations for Inference
|
||||
- Keep `AutencoderKLWan` in `torch.float32` for better decoding quality.
|
||||
- `num_frames` should satisfy the following constraint: `(num_frames - 1) % 4 == 0`
|
||||
- For smaller resolution videos, try lower values of `shift` (between `2.0` to `5.0`) in the [Scheduler](https://huggingface.co/docs/diffusers/main/en/api/schedulers/flow_match_euler_discrete#diffusers.FlowMatchEulerDiscreteScheduler.shift). For larger resolution videos, try higher values (between `7.0` and `12.0`). The default value is `3.0` for Wan.
|
||||
|
||||
## WanPipeline
|
||||
|
||||
[[autodoc]] WanPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## WanImageToVideoPipeline
|
||||
|
||||
[[autodoc]] WanImageToVideoPipeline
|
||||
- all
|
||||
- __call__
|
||||
|
||||
## WanPipelineOutput
|
||||
|
||||
[[autodoc]] pipelines.wan.pipeline_output.WanPipelineOutput
|
||||
@@ -12,6 +12,10 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Würstchen
|
||||
|
||||
<div class="flex flex-wrap space-x-1">
|
||||
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
|
||||
</div>
|
||||
|
||||
<img src="https://github.com/dome272/Wuerstchen/assets/61938694/0617c863-165a-43ee-9303-2a17299a0cf9">
|
||||
|
||||
[Wuerstchen: An Efficient Architecture for Large-Scale Text-to-Image Diffusion Models](https://huggingface.co/papers/2306.00637) is by Pablo Pernias, Dominic Rampas, Mats L. Richter and Christopher Pal and Marc Aubreville.
|
||||
|
||||
@@ -31,6 +31,11 @@ Learn how to quantize models in the [Quantization](../quantization/overview) gui
|
||||
## GGUFQuantizationConfig
|
||||
|
||||
[[autodoc]] GGUFQuantizationConfig
|
||||
|
||||
## QuantoConfig
|
||||
|
||||
[[autodoc]] QuantoConfig
|
||||
|
||||
## TorchAoConfig
|
||||
|
||||
[[autodoc]] TorchAoConfig
|
||||
|
||||
19
docs/source/en/api/schedulers/ddim_cogvideox.md
Normal file
19
docs/source/en/api/schedulers/ddim_cogvideox.md
Normal file
@@ -0,0 +1,19 @@
|
||||
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# CogVideoXDDIMScheduler
|
||||
|
||||
`CogVideoXDDIMScheduler` is based on [Denoising Diffusion Implicit Models](https://huggingface.co/papers/2010.02502), specifically for CogVideoX models.
|
||||
|
||||
## CogVideoXDDIMScheduler
|
||||
|
||||
[[autodoc]] CogVideoXDDIMScheduler
|
||||
@@ -0,0 +1,19 @@
|
||||
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# CogVideoXDPMScheduler
|
||||
|
||||
`CogVideoXDPMScheduler` is based on [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095), specifically for CogVideoX models.
|
||||
|
||||
## CogVideoXDPMScheduler
|
||||
|
||||
[[autodoc]] CogVideoXDPMScheduler
|
||||
@@ -41,3 +41,11 @@ Utility and helper functions for working with 🤗 Diffusers.
|
||||
## randn_tensor
|
||||
|
||||
[[autodoc]] utils.torch_utils.randn_tensor
|
||||
|
||||
## apply_layerwise_casting
|
||||
|
||||
[[autodoc]] hooks.layerwise_casting.apply_layerwise_casting
|
||||
|
||||
## apply_group_offloading
|
||||
|
||||
[[autodoc]] hooks.group_offloading.apply_group_offloading
|
||||
|
||||
@@ -83,4 +83,8 @@ Happy exploring, and thank you for being part of the Diffusers community!
|
||||
<td><a href="https://github.com/suzukimain/auto_diffusers"> Model Search </a></td>
|
||||
<td>Search models on Civitai and Hugging Face</td>
|
||||
</tr>
|
||||
<tr style="border-top: 2px solid black">
|
||||
<td><a href="https://github.com/beinsezii/skrample"> Skrample </a></td>
|
||||
<td>Fully modular scheduler functions with 1st class diffusers integration.</td>
|
||||
</tr>
|
||||
</table>
|
||||
|
||||
@@ -16,6 +16,11 @@ specific language governing permissions and limitations under the License.
|
||||
<img src="https://colab.research.google.com/assets/colab-badge.svg" alt="Open In Colab"/>
|
||||
</a>
|
||||
|
||||
> [!TIP]
|
||||
> This document has now grown outdated given the emergence of existing evaluation frameworks for diffusion models for image generation. Please check
|
||||
> out works like [HEIM](https://crfm.stanford.edu/helm/heim/latest/), [T2I-Compbench](https://arxiv.org/abs/2307.06350),
|
||||
> [GenEval](https://arxiv.org/abs/2310.11513).
|
||||
|
||||
Evaluation of generative models like [Stable Diffusion](https://huggingface.co/docs/diffusers/stable_diffusion) is subjective in nature. But as practitioners and researchers, we often have to make careful choices amongst many different possibilities. So, when working with different generative models (like GANs, Diffusion, etc.), how do we choose one over the other?
|
||||
|
||||
Qualitative evaluation of such models can be error-prone and might incorrectly influence a decision.
|
||||
|
||||
9
docs/source/en/hybrid_inference/api_reference.md
Normal file
9
docs/source/en/hybrid_inference/api_reference.md
Normal file
@@ -0,0 +1,9 @@
|
||||
# Hybrid Inference API Reference
|
||||
|
||||
## Remote Decode
|
||||
|
||||
[[autodoc]] utils.remote_utils.remote_decode
|
||||
|
||||
## Remote Encode
|
||||
|
||||
[[autodoc]] utils.remote_utils.remote_encode
|
||||
60
docs/source/en/hybrid_inference/overview.md
Normal file
60
docs/source/en/hybrid_inference/overview.md
Normal file
@@ -0,0 +1,60 @@
|
||||
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Hybrid Inference
|
||||
|
||||
**Empowering local AI builders with Hybrid Inference**
|
||||
|
||||
|
||||
> [!TIP]
|
||||
> Hybrid Inference is an [experimental feature](https://huggingface.co/blog/remote_vae).
|
||||
> Feedback can be provided [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
|
||||
|
||||
|
||||
|
||||
## Why use Hybrid Inference?
|
||||
|
||||
Hybrid Inference offers a fast and simple way to offload local generation requirements.
|
||||
|
||||
- 🚀 **Reduced Requirements:** Access powerful models without expensive hardware.
|
||||
- 💎 **Without Compromise:** Achieve the highest quality without sacrificing performance.
|
||||
- 💰 **Cost Effective:** It's free! 🤑
|
||||
- 🎯 **Diverse Use Cases:** Fully compatible with Diffusers 🧨 and the wider community.
|
||||
- 🔧 **Developer-Friendly:** Simple requests, fast responses.
|
||||
|
||||
---
|
||||
|
||||
## Available Models
|
||||
|
||||
* **VAE Decode 🖼️:** Quickly decode latent representations into high-quality images without compromising performance or workflow speed.
|
||||
* **VAE Encode 🔢:** Efficiently encode images into latent representations for generation and training.
|
||||
* **Text Encoders 📃 (coming soon):** Compute text embeddings for your prompts quickly and accurately, ensuring a smooth and high-quality workflow.
|
||||
|
||||
---
|
||||
|
||||
## Integrations
|
||||
|
||||
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
|
||||
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.
|
||||
|
||||
## Changelog
|
||||
|
||||
- March 10 2025: Added VAE encode
|
||||
- March 2 2025: Initial release with VAE decoding
|
||||
|
||||
## Contents
|
||||
|
||||
The documentation is organized into three sections:
|
||||
|
||||
* **VAE Decode** Learn the basics of how to use VAE Decode with Hybrid Inference.
|
||||
* **VAE Encode** Learn the basics of how to use VAE Encode with Hybrid Inference.
|
||||
* **API Reference** Dive into task-specific settings and parameters.
|
||||
345
docs/source/en/hybrid_inference/vae_decode.md
Normal file
345
docs/source/en/hybrid_inference/vae_decode.md
Normal file
@@ -0,0 +1,345 @@
|
||||
# Getting Started: VAE Decode with Hybrid Inference
|
||||
|
||||
VAE decode is an essential component of diffusion models - turning latent representations into images or videos.
|
||||
|
||||
## Memory
|
||||
|
||||
These tables demonstrate the VRAM requirements for VAE decode with SD v1 and SD XL on different GPUs.
|
||||
|
||||
For the majority of these GPUs the memory usage % dictates other models (text encoders, UNet/Transformer) must be offloaded, or tiled decoding has to be used which increases time taken and impacts quality.
|
||||
|
||||
<details><summary>SD v1.5</summary>
|
||||
|
||||
| GPU | Resolution | Time (seconds) | Memory (%) | Tiled Time (secs) | Tiled Memory (%) |
|
||||
| --- | --- | --- | --- | --- | --- |
|
||||
| NVIDIA GeForce RTX 4090 | 512x512 | 0.031 | 5.60% | 0.031 (0%) | 5.60% |
|
||||
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.148 | 20.00% | 0.301 (+103%) | 5.60% |
|
||||
| NVIDIA GeForce RTX 4080 | 512x512 | 0.05 | 8.40% | 0.050 (0%) | 8.40% |
|
||||
| NVIDIA GeForce RTX 4080 | 1024x1024 | 0.224 | 30.00% | 0.356 (+59%) | 8.40% |
|
||||
| NVIDIA GeForce RTX 4070 Ti | 512x512 | 0.066 | 11.30% | 0.066 (0%) | 11.30% |
|
||||
| NVIDIA GeForce RTX 4070 Ti | 1024x1024 | 0.284 | 40.50% | 0.454 (+60%) | 11.40% |
|
||||
| NVIDIA GeForce RTX 3090 | 512x512 | 0.062 | 5.20% | 0.062 (0%) | 5.20% |
|
||||
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.253 | 18.50% | 0.464 (+83%) | 5.20% |
|
||||
| NVIDIA GeForce RTX 3080 | 512x512 | 0.07 | 12.80% | 0.070 (0%) | 12.80% |
|
||||
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.286 | 45.30% | 0.466 (+63%) | 12.90% |
|
||||
| NVIDIA GeForce RTX 3070 | 512x512 | 0.102 | 15.90% | 0.102 (0%) | 15.90% |
|
||||
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.421 | 56.30% | 0.746 (+77%) | 16.00% |
|
||||
|
||||
</details>
|
||||
|
||||
<details><summary>SDXL</summary>
|
||||
|
||||
| GPU | Resolution | Time (seconds) | Memory Consumed (%) | Tiled Time (seconds) | Tiled Memory (%) |
|
||||
| --- | --- | --- | --- | --- | --- |
|
||||
| NVIDIA GeForce RTX 4090 | 512x512 | 0.057 | 10.00% | 0.057 (0%) | 10.00% |
|
||||
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.256 | 35.50% | 0.257 (+0.4%) | 35.50% |
|
||||
| NVIDIA GeForce RTX 4080 | 512x512 | 0.092 | 15.00% | 0.092 (0%) | 15.00% |
|
||||
| NVIDIA GeForce RTX 4080 | 1024x1024 | 0.406 | 53.30% | 0.406 (0%) | 53.30% |
|
||||
| NVIDIA GeForce RTX 4070 Ti | 512x512 | 0.121 | 20.20% | 0.120 (-0.8%) | 20.20% |
|
||||
| NVIDIA GeForce RTX 4070 Ti | 1024x1024 | 0.519 | 72.00% | 0.519 (0%) | 72.00% |
|
||||
| NVIDIA GeForce RTX 3090 | 512x512 | 0.107 | 10.50% | 0.107 (0%) | 10.50% |
|
||||
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.459 | 38.00% | 0.460 (+0.2%) | 38.00% |
|
||||
| NVIDIA GeForce RTX 3080 | 512x512 | 0.121 | 25.60% | 0.121 (0%) | 25.60% |
|
||||
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.524 | 93.00% | 0.524 (0%) | 93.00% |
|
||||
| NVIDIA GeForce RTX 3070 | 512x512 | 0.183 | 31.80% | 0.183 (0%) | 31.80% |
|
||||
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.794 | 96.40% | 0.794 (0%) | 96.40% |
|
||||
|
||||
</details>
|
||||
|
||||
## Available VAEs
|
||||
|
||||
| | **Endpoint** | **Model** |
|
||||
|:-:|:-----------:|:--------:|
|
||||
| **Stable Diffusion v1** | [https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud](https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud) | [`stabilityai/sd-vae-ft-mse`](https://hf.co/stabilityai/sd-vae-ft-mse) |
|
||||
| **Stable Diffusion XL** | [https://x2dmsqunjd6k9prw.us-east-1.aws.endpoints.huggingface.cloud](https://x2dmsqunjd6k9prw.us-east-1.aws.endpoints.huggingface.cloud) | [`madebyollin/sdxl-vae-fp16-fix`](https://hf.co/madebyollin/sdxl-vae-fp16-fix) |
|
||||
| **Flux** | [https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud](https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud) | [`black-forest-labs/FLUX.1-schnell`](https://hf.co/black-forest-labs/FLUX.1-schnell) |
|
||||
| **HunyuanVideo** | [https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud](https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud) | [`hunyuanvideo-community/HunyuanVideo`](https://hf.co/hunyuanvideo-community/HunyuanVideo) |
|
||||
|
||||
|
||||
> [!TIP]
|
||||
> Model support can be requested [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
|
||||
|
||||
|
||||
## Code
|
||||
|
||||
> [!TIP]
|
||||
> Install `diffusers` from `main` to run the code: `pip install git+https://github.com/huggingface/diffusers@main`
|
||||
|
||||
|
||||
A helper method simplifies interacting with Hybrid Inference.
|
||||
|
||||
```python
|
||||
from diffusers.utils.remote_utils import remote_decode
|
||||
```
|
||||
|
||||
### Basic example
|
||||
|
||||
Here, we show how to use the remote VAE on random tensors.
|
||||
|
||||
<details><summary>Code</summary>
|
||||
|
||||
```python
|
||||
image = remote_decode(
|
||||
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
tensor=torch.randn([1, 4, 64, 64], dtype=torch.float16),
|
||||
scaling_factor=0.18215,
|
||||
)
|
||||
```
|
||||
|
||||
</details>
|
||||
|
||||
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
|
||||
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/output.png"/>
|
||||
</figure>
|
||||
|
||||
Usage for Flux is slightly different. Flux latents are packed so we need to send the `height` and `width`.
|
||||
|
||||
<details><summary>Code</summary>
|
||||
|
||||
```python
|
||||
image = remote_decode(
|
||||
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
tensor=torch.randn([1, 4096, 64], dtype=torch.float16),
|
||||
height=1024,
|
||||
width=1024,
|
||||
scaling_factor=0.3611,
|
||||
shift_factor=0.1159,
|
||||
)
|
||||
```
|
||||
|
||||
</details>
|
||||
|
||||
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
|
||||
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/flux_random_latent.png"/>
|
||||
</figure>
|
||||
|
||||
Finally, an example for HunyuanVideo.
|
||||
|
||||
<details><summary>Code</summary>
|
||||
|
||||
```python
|
||||
video = remote_decode(
|
||||
endpoint="https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
tensor=torch.randn([1, 16, 3, 40, 64], dtype=torch.float16),
|
||||
output_type="mp4",
|
||||
)
|
||||
with open("video.mp4", "wb") as f:
|
||||
f.write(video)
|
||||
```
|
||||
|
||||
</details>
|
||||
|
||||
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
|
||||
<video
|
||||
alt="queue.mp4"
|
||||
autoplay loop autobuffer muted playsinline
|
||||
>
|
||||
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/video_1.mp4" type="video/mp4">
|
||||
</video>
|
||||
</figure>
|
||||
|
||||
|
||||
### Generation
|
||||
|
||||
But we want to use the VAE on an actual pipeline to get an actual image, not random noise. The example below shows how to do it with SD v1.5.
|
||||
|
||||
<details><summary>Code</summary>
|
||||
|
||||
```python
|
||||
from diffusers import StableDiffusionPipeline
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained(
|
||||
"stable-diffusion-v1-5/stable-diffusion-v1-5",
|
||||
torch_dtype=torch.float16,
|
||||
variant="fp16",
|
||||
vae=None,
|
||||
).to("cuda")
|
||||
|
||||
prompt = "Strawberry ice cream, in a stylish modern glass, coconut, splashing milk cream and honey, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious"
|
||||
|
||||
latent = pipe(
|
||||
prompt=prompt,
|
||||
output_type="latent",
|
||||
).images
|
||||
image = remote_decode(
|
||||
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
tensor=latent,
|
||||
scaling_factor=0.18215,
|
||||
)
|
||||
image.save("test.jpg")
|
||||
```
|
||||
|
||||
</details>
|
||||
|
||||
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
|
||||
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/test.jpg"/>
|
||||
</figure>
|
||||
|
||||
Here’s another example with Flux.
|
||||
|
||||
<details><summary>Code</summary>
|
||||
|
||||
```python
|
||||
from diffusers import FluxPipeline
|
||||
|
||||
pipe = FluxPipeline.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-schnell",
|
||||
torch_dtype=torch.bfloat16,
|
||||
vae=None,
|
||||
).to("cuda")
|
||||
|
||||
prompt = "Strawberry ice cream, in a stylish modern glass, coconut, splashing milk cream and honey, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious"
|
||||
|
||||
latent = pipe(
|
||||
prompt=prompt,
|
||||
guidance_scale=0.0,
|
||||
num_inference_steps=4,
|
||||
output_type="latent",
|
||||
).images
|
||||
image = remote_decode(
|
||||
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
tensor=latent,
|
||||
height=1024,
|
||||
width=1024,
|
||||
scaling_factor=0.3611,
|
||||
shift_factor=0.1159,
|
||||
)
|
||||
image.save("test.jpg")
|
||||
```
|
||||
|
||||
</details>
|
||||
|
||||
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
|
||||
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/test_1.jpg"/>
|
||||
</figure>
|
||||
|
||||
Here’s an example with HunyuanVideo.
|
||||
|
||||
<details><summary>Code</summary>
|
||||
|
||||
```python
|
||||
from diffusers import HunyuanVideoPipeline, HunyuanVideoTransformer3DModel
|
||||
|
||||
model_id = "hunyuanvideo-community/HunyuanVideo"
|
||||
transformer = HunyuanVideoTransformer3DModel.from_pretrained(
|
||||
model_id, subfolder="transformer", torch_dtype=torch.bfloat16
|
||||
)
|
||||
pipe = HunyuanVideoPipeline.from_pretrained(
|
||||
model_id, transformer=transformer, vae=None, torch_dtype=torch.float16
|
||||
).to("cuda")
|
||||
|
||||
latent = pipe(
|
||||
prompt="A cat walks on the grass, realistic",
|
||||
height=320,
|
||||
width=512,
|
||||
num_frames=61,
|
||||
num_inference_steps=30,
|
||||
output_type="latent",
|
||||
).frames
|
||||
|
||||
video = remote_decode(
|
||||
endpoint="https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
tensor=latent,
|
||||
output_type="mp4",
|
||||
)
|
||||
|
||||
if isinstance(video, bytes):
|
||||
with open("video.mp4", "wb") as f:
|
||||
f.write(video)
|
||||
```
|
||||
|
||||
</details>
|
||||
|
||||
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
|
||||
<video
|
||||
alt="queue.mp4"
|
||||
autoplay loop autobuffer muted playsinline
|
||||
>
|
||||
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/video.mp4" type="video/mp4">
|
||||
</video>
|
||||
</figure>
|
||||
|
||||
|
||||
### Queueing
|
||||
|
||||
One of the great benefits of using a remote VAE is that we can queue multiple generation requests. While the current latent is being processed for decoding, we can already queue another one. This helps improve concurrency.
|
||||
|
||||
|
||||
<details><summary>Code</summary>
|
||||
|
||||
```python
|
||||
import queue
|
||||
import threading
|
||||
from IPython.display import display
|
||||
from diffusers import StableDiffusionPipeline
|
||||
|
||||
def decode_worker(q: queue.Queue):
|
||||
while True:
|
||||
item = q.get()
|
||||
if item is None:
|
||||
break
|
||||
image = remote_decode(
|
||||
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
tensor=item,
|
||||
scaling_factor=0.18215,
|
||||
)
|
||||
display(image)
|
||||
q.task_done()
|
||||
|
||||
q = queue.Queue()
|
||||
thread = threading.Thread(target=decode_worker, args=(q,), daemon=True)
|
||||
thread.start()
|
||||
|
||||
def decode(latent: torch.Tensor):
|
||||
q.put(latent)
|
||||
|
||||
prompts = [
|
||||
"Blueberry ice cream, in a stylish modern glass , ice cubes, nuts, mint leaves, splashing milk cream, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious",
|
||||
"Lemonade in a glass, mint leaves, in an aqua and white background, flowers, ice cubes, halo, fluid motion, dynamic movement, soft lighting, digital painting, rule of thirds composition, Art by Greg rutkowski, Coby whitmore",
|
||||
"Comic book art, beautiful, vintage, pastel neon colors, extremely detailed pupils, delicate features, light on face, slight smile, Artgerm, Mary Blair, Edmund Dulac, long dark locks, bangs, glowing, fashionable style, fairytale ambience, hot pink.",
|
||||
"Masterpiece, vanilla cone ice cream garnished with chocolate syrup, crushed nuts, choco flakes, in a brown background, gold, cinematic lighting, Art by WLOP",
|
||||
"A bowl of milk, falling cornflakes, berries, blueberries, in a white background, soft lighting, intricate details, rule of thirds, octane render, volumetric lighting",
|
||||
"Cold Coffee with cream, crushed almonds, in a glass, choco flakes, ice cubes, wet, in a wooden background, cinematic lighting, hyper realistic painting, art by Carne Griffiths, octane render, volumetric lighting, fluid motion, dynamic movement, muted colors,",
|
||||
]
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained(
|
||||
"Lykon/dreamshaper-8",
|
||||
torch_dtype=torch.float16,
|
||||
vae=None,
|
||||
).to("cuda")
|
||||
|
||||
pipe.unet = pipe.unet.to(memory_format=torch.channels_last)
|
||||
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
|
||||
|
||||
_ = pipe(
|
||||
prompt=prompts[0],
|
||||
output_type="latent",
|
||||
)
|
||||
|
||||
for prompt in prompts:
|
||||
latent = pipe(
|
||||
prompt=prompt,
|
||||
output_type="latent",
|
||||
).images
|
||||
decode(latent)
|
||||
|
||||
q.put(None)
|
||||
thread.join()
|
||||
```
|
||||
|
||||
</details>
|
||||
|
||||
|
||||
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
|
||||
<video
|
||||
alt="queue.mp4"
|
||||
autoplay loop autobuffer muted playsinline
|
||||
>
|
||||
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/queue.mp4" type="video/mp4">
|
||||
</video>
|
||||
</figure>
|
||||
|
||||
## Integrations
|
||||
|
||||
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
|
||||
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.
|
||||
183
docs/source/en/hybrid_inference/vae_encode.md
Normal file
183
docs/source/en/hybrid_inference/vae_encode.md
Normal file
@@ -0,0 +1,183 @@
|
||||
# Getting Started: VAE Encode with Hybrid Inference
|
||||
|
||||
VAE encode is used for training, image-to-image and image-to-video - turning into images or videos into latent representations.
|
||||
|
||||
## Memory
|
||||
|
||||
These tables demonstrate the VRAM requirements for VAE encode with SD v1 and SD XL on different GPUs.
|
||||
|
||||
For the majority of these GPUs the memory usage % dictates other models (text encoders, UNet/Transformer) must be offloaded, or tiled encoding has to be used which increases time taken and impacts quality.
|
||||
|
||||
<details><summary>SD v1.5</summary>
|
||||
|
||||
| GPU | Resolution | Time (seconds) | Memory (%) | Tiled Time (secs) | Tiled Memory (%) |
|
||||
|:------------------------------|:-------------|-----------------:|-------------:|--------------------:|-------------------:|
|
||||
| NVIDIA GeForce RTX 4090 | 512x512 | 0.015 | 3.51901 | 0.015 | 3.51901 |
|
||||
| NVIDIA GeForce RTX 4090 | 256x256 | 0.004 | 1.3154 | 0.005 | 1.3154 |
|
||||
| NVIDIA GeForce RTX 4090 | 2048x2048 | 0.402 | 47.1852 | 0.496 | 3.51901 |
|
||||
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.078 | 12.2658 | 0.094 | 3.51901 |
|
||||
| NVIDIA GeForce RTX 4080 SUPER | 512x512 | 0.023 | 5.30105 | 0.023 | 5.30105 |
|
||||
| NVIDIA GeForce RTX 4080 SUPER | 256x256 | 0.006 | 1.98152 | 0.006 | 1.98152 |
|
||||
| NVIDIA GeForce RTX 4080 SUPER | 2048x2048 | 0.574 | 71.08 | 0.656 | 5.30105 |
|
||||
| NVIDIA GeForce RTX 4080 SUPER | 1024x1024 | 0.111 | 18.4772 | 0.14 | 5.30105 |
|
||||
| NVIDIA GeForce RTX 3090 | 512x512 | 0.032 | 3.52782 | 0.032 | 3.52782 |
|
||||
| NVIDIA GeForce RTX 3090 | 256x256 | 0.01 | 1.31869 | 0.009 | 1.31869 |
|
||||
| NVIDIA GeForce RTX 3090 | 2048x2048 | 0.742 | 47.3033 | 0.954 | 3.52782 |
|
||||
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.136 | 12.2965 | 0.207 | 3.52782 |
|
||||
| NVIDIA GeForce RTX 3080 | 512x512 | 0.036 | 8.51761 | 0.036 | 8.51761 |
|
||||
| NVIDIA GeForce RTX 3080 | 256x256 | 0.01 | 3.18387 | 0.01 | 3.18387 |
|
||||
| NVIDIA GeForce RTX 3080 | 2048x2048 | 0.863 | 86.7424 | 1.191 | 8.51761 |
|
||||
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.157 | 29.6888 | 0.227 | 8.51761 |
|
||||
| NVIDIA GeForce RTX 3070 | 512x512 | 0.051 | 10.6941 | 0.051 | 10.6941 |
|
||||
| NVIDIA GeForce RTX 3070 | 256x256 | 0.015 | 3.99743 | 0.015 | 3.99743 |
|
||||
| NVIDIA GeForce RTX 3070 | 2048x2048 | 1.217 | 96.054 | 1.482 | 10.6941 |
|
||||
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.223 | 37.2751 | 0.327 | 10.6941 |
|
||||
|
||||
|
||||
</details>
|
||||
|
||||
<details><summary>SDXL</summary>
|
||||
|
||||
| GPU | Resolution | Time (seconds) | Memory Consumed (%) | Tiled Time (seconds) | Tiled Memory (%) |
|
||||
|:------------------------------|:-------------|-----------------:|----------------------:|-----------------------:|-------------------:|
|
||||
| NVIDIA GeForce RTX 4090 | 512x512 | 0.029 | 4.95707 | 0.029 | 4.95707 |
|
||||
| NVIDIA GeForce RTX 4090 | 256x256 | 0.007 | 2.29666 | 0.007 | 2.29666 |
|
||||
| NVIDIA GeForce RTX 4090 | 2048x2048 | 0.873 | 66.3452 | 0.863 | 15.5649 |
|
||||
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.142 | 15.5479 | 0.143 | 15.5479 |
|
||||
| NVIDIA GeForce RTX 4080 SUPER | 512x512 | 0.044 | 7.46735 | 0.044 | 7.46735 |
|
||||
| NVIDIA GeForce RTX 4080 SUPER | 256x256 | 0.01 | 3.4597 | 0.01 | 3.4597 |
|
||||
| NVIDIA GeForce RTX 4080 SUPER | 2048x2048 | 1.317 | 87.1615 | 1.291 | 23.447 |
|
||||
| NVIDIA GeForce RTX 4080 SUPER | 1024x1024 | 0.213 | 23.4215 | 0.214 | 23.4215 |
|
||||
| NVIDIA GeForce RTX 3090 | 512x512 | 0.058 | 5.65638 | 0.058 | 5.65638 |
|
||||
| NVIDIA GeForce RTX 3090 | 256x256 | 0.016 | 2.45081 | 0.016 | 2.45081 |
|
||||
| NVIDIA GeForce RTX 3090 | 2048x2048 | 1.755 | 77.8239 | 1.614 | 18.4193 |
|
||||
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.265 | 18.4023 | 0.265 | 18.4023 |
|
||||
| NVIDIA GeForce RTX 3080 | 512x512 | 0.064 | 13.6568 | 0.064 | 13.6568 |
|
||||
| NVIDIA GeForce RTX 3080 | 256x256 | 0.018 | 5.91728 | 0.018 | 5.91728 |
|
||||
| NVIDIA GeForce RTX 3080 | 2048x2048 | OOM | OOM | 1.866 | 44.4717 |
|
||||
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.302 | 44.4308 | 0.302 | 44.4308 |
|
||||
| NVIDIA GeForce RTX 3070 | 512x512 | 0.093 | 17.1465 | 0.093 | 17.1465 |
|
||||
| NVIDIA GeForce RTX 3070 | 256x256 | 0.025 | 7.42931 | 0.026 | 7.42931 |
|
||||
| NVIDIA GeForce RTX 3070 | 2048x2048 | OOM | OOM | 2.674 | 55.8355 |
|
||||
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.443 | 55.7841 | 0.443 | 55.7841 |
|
||||
|
||||
</details>
|
||||
|
||||
## Available VAEs
|
||||
|
||||
| | **Endpoint** | **Model** |
|
||||
|:-:|:-----------:|:--------:|
|
||||
| **Stable Diffusion v1** | [https://qc6479g0aac6qwy9.us-east-1.aws.endpoints.huggingface.cloud](https://qc6479g0aac6qwy9.us-east-1.aws.endpoints.huggingface.cloud) | [`stabilityai/sd-vae-ft-mse`](https://hf.co/stabilityai/sd-vae-ft-mse) |
|
||||
| **Stable Diffusion XL** | [https://xjqqhmyn62rog84g.us-east-1.aws.endpoints.huggingface.cloud](https://xjqqhmyn62rog84g.us-east-1.aws.endpoints.huggingface.cloud) | [`madebyollin/sdxl-vae-fp16-fix`](https://hf.co/madebyollin/sdxl-vae-fp16-fix) |
|
||||
| **Flux** | [https://ptccx55jz97f9zgo.us-east-1.aws.endpoints.huggingface.cloud](https://ptccx55jz97f9zgo.us-east-1.aws.endpoints.huggingface.cloud) | [`black-forest-labs/FLUX.1-schnell`](https://hf.co/black-forest-labs/FLUX.1-schnell) |
|
||||
|
||||
|
||||
> [!TIP]
|
||||
> Model support can be requested [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
|
||||
|
||||
|
||||
## Code
|
||||
|
||||
> [!TIP]
|
||||
> Install `diffusers` from `main` to run the code: `pip install git+https://github.com/huggingface/diffusers@main`
|
||||
|
||||
|
||||
A helper method simplifies interacting with Hybrid Inference.
|
||||
|
||||
```python
|
||||
from diffusers.utils.remote_utils import remote_encode
|
||||
```
|
||||
|
||||
### Basic example
|
||||
|
||||
Let's encode an image, then decode it to demonstrate.
|
||||
|
||||
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
|
||||
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"/>
|
||||
</figure>
|
||||
|
||||
<details><summary>Code</summary>
|
||||
|
||||
```python
|
||||
from diffusers.utils import load_image
|
||||
from diffusers.utils.remote_utils import remote_decode
|
||||
|
||||
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg?download=true")
|
||||
|
||||
latent = remote_encode(
|
||||
endpoint="https://ptccx55jz97f9zgo.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
scaling_factor=0.3611,
|
||||
shift_factor=0.1159,
|
||||
)
|
||||
|
||||
decoded = remote_decode(
|
||||
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
tensor=latent,
|
||||
scaling_factor=0.3611,
|
||||
shift_factor=0.1159,
|
||||
)
|
||||
```
|
||||
|
||||
</details>
|
||||
|
||||
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
|
||||
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/decoded.png"/>
|
||||
</figure>
|
||||
|
||||
|
||||
### Generation
|
||||
|
||||
Now let's look at a generation example, we'll encode the image, generate then remotely decode too!
|
||||
|
||||
<details><summary>Code</summary>
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import StableDiffusionImg2ImgPipeline
|
||||
from diffusers.utils import load_image
|
||||
from diffusers.utils.remote_utils import remote_decode, remote_encode
|
||||
|
||||
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
|
||||
"stable-diffusion-v1-5/stable-diffusion-v1-5",
|
||||
torch_dtype=torch.float16,
|
||||
variant="fp16",
|
||||
vae=None,
|
||||
).to("cuda")
|
||||
|
||||
init_image = load_image(
|
||||
"https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
|
||||
)
|
||||
init_image = init_image.resize((768, 512))
|
||||
|
||||
init_latent = remote_encode(
|
||||
endpoint="https://qc6479g0aac6qwy9.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
image=init_image,
|
||||
scaling_factor=0.18215,
|
||||
)
|
||||
|
||||
prompt = "A fantasy landscape, trending on artstation"
|
||||
latent = pipe(
|
||||
prompt=prompt,
|
||||
image=init_latent,
|
||||
strength=0.75,
|
||||
output_type="latent",
|
||||
).images
|
||||
|
||||
image = remote_decode(
|
||||
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
|
||||
tensor=latent,
|
||||
scaling_factor=0.18215,
|
||||
)
|
||||
image.save("fantasy_landscape.jpg")
|
||||
```
|
||||
|
||||
</details>
|
||||
|
||||
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
|
||||
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/fantasy_landscape.png"/>
|
||||
</figure>
|
||||
|
||||
## Integrations
|
||||
|
||||
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
|
||||
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.
|
||||
@@ -23,32 +23,60 @@ You should install 🤗 Diffusers in a [virtual environment](https://docs.python
|
||||
If you're unfamiliar with Python virtual environments, take a look at this [guide](https://packaging.python.org/guides/installing-using-pip-and-virtual-environments/).
|
||||
A virtual environment makes it easier to manage different projects and avoid compatibility issues between dependencies.
|
||||
|
||||
Start by creating a virtual environment in your project directory:
|
||||
Create a virtual environment with Python or [uv](https://docs.astral.sh/uv/) (refer to [Installation](https://docs.astral.sh/uv/getting-started/installation/) for installation instructions), a fast Rust-based Python package and project manager.
|
||||
|
||||
<hfoptions id="install">
|
||||
<hfoption id="uv">
|
||||
|
||||
```bash
|
||||
python -m venv .env
|
||||
uv venv my-env
|
||||
source my-env/bin/activate
|
||||
```
|
||||
|
||||
Activate the virtual environment:
|
||||
</hfoption>
|
||||
<hfoption id="Python">
|
||||
|
||||
```bash
|
||||
source .env/bin/activate
|
||||
python -m venv my-env
|
||||
source my-env/bin/activate
|
||||
```
|
||||
|
||||
You should also install 🤗 Transformers because 🤗 Diffusers relies on its models:
|
||||
</hfoption>
|
||||
</hfoptions>
|
||||
|
||||
You should also install 🤗 Transformers because 🤗 Diffusers relies on its models.
|
||||
|
||||
|
||||
<frameworkcontent>
|
||||
<pt>
|
||||
Note - PyTorch only supports Python 3.8 - 3.11 on Windows.
|
||||
|
||||
PyTorch only supports Python 3.8 - 3.11 on Windows. Install Diffusers with uv.
|
||||
|
||||
```bash
|
||||
uv install diffusers["torch"] transformers
|
||||
```
|
||||
|
||||
You can also install Diffusers with pip.
|
||||
|
||||
```bash
|
||||
pip install diffusers["torch"] transformers
|
||||
```
|
||||
|
||||
</pt>
|
||||
<jax>
|
||||
|
||||
Install Diffusers with uv.
|
||||
|
||||
```bash
|
||||
uv pip install diffusers["flax"] transformers
|
||||
```
|
||||
|
||||
You can also install Diffusers with pip.
|
||||
|
||||
```bash
|
||||
pip install diffusers["flax"] transformers
|
||||
```
|
||||
|
||||
</jax>
|
||||
</frameworkcontent>
|
||||
|
||||
@@ -133,10 +161,10 @@ Your Python environment will find the `main` version of 🤗 Diffusers on the ne
|
||||
|
||||
Model weights and files are downloaded from the Hub to a cache which is usually your home directory. You can change the cache location by specifying the `HF_HOME` or `HUGGINFACE_HUB_CACHE` environment variables or configuring the `cache_dir` parameter in methods like [`~DiffusionPipeline.from_pretrained`].
|
||||
|
||||
Cached files allow you to run 🤗 Diffusers offline. To prevent 🤗 Diffusers from connecting to the internet, set the `HF_HUB_OFFLINE` environment variable to `True` and 🤗 Diffusers will only load previously downloaded files in the cache.
|
||||
Cached files allow you to run 🤗 Diffusers offline. To prevent 🤗 Diffusers from connecting to the internet, set the `HF_HUB_OFFLINE` environment variable to `1` and 🤗 Diffusers will only load previously downloaded files in the cache.
|
||||
|
||||
```shell
|
||||
export HF_HUB_OFFLINE=True
|
||||
export HF_HUB_OFFLINE=1
|
||||
```
|
||||
|
||||
For more details about managing and cleaning the cache, take a look at the [caching](https://huggingface.co/docs/huggingface_hub/guides/manage-cache) guide.
|
||||
@@ -151,14 +179,16 @@ Telemetry is only sent when loading models and pipelines from the Hub,
|
||||
and it is not collected if you're loading local files.
|
||||
|
||||
We understand that not everyone wants to share additional information,and we respect your privacy.
|
||||
You can disable telemetry collection by setting the `DISABLE_TELEMETRY` environment variable from your terminal:
|
||||
You can disable telemetry collection by setting the `HF_HUB_DISABLE_TELEMETRY` environment variable from your terminal:
|
||||
|
||||
On Linux/MacOS:
|
||||
|
||||
```bash
|
||||
export DISABLE_TELEMETRY=YES
|
||||
export HF_HUB_DISABLE_TELEMETRY=1
|
||||
```
|
||||
|
||||
On Windows:
|
||||
|
||||
```bash
|
||||
set DISABLE_TELEMETRY=YES
|
||||
set HF_HUB_DISABLE_TELEMETRY=1
|
||||
```
|
||||
|
||||
@@ -158,6 +158,107 @@ In order to properly offload models after they're called, it is required to run
|
||||
|
||||
</Tip>
|
||||
|
||||
## Group offloading
|
||||
|
||||
Group offloading is the middle ground between sequential and model offloading. It works by offloading groups of internal layers (either `torch.nn.ModuleList` or `torch.nn.Sequential`), which uses less memory than model-level offloading. It is also faster than sequential-level offloading because the number of device synchronizations is reduced.
|
||||
|
||||
To enable group offloading, call the [`~ModelMixin.enable_group_offload`] method on the model if it is a Diffusers model implementation. For any other model implementation, use [`~hooks.group_offloading.apply_group_offloading`]:
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import CogVideoXPipeline
|
||||
from diffusers.hooks import apply_group_offloading
|
||||
from diffusers.utils import export_to_video
|
||||
|
||||
# Load the pipeline
|
||||
onload_device = torch.device("cuda")
|
||||
offload_device = torch.device("cpu")
|
||||
pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
|
||||
|
||||
# We can utilize the enable_group_offload method for Diffusers model implementations
|
||||
pipe.transformer.enable_group_offload(onload_device=onload_device, offload_device=offload_device, offload_type="leaf_level", use_stream=True)
|
||||
|
||||
# Uncomment the following to also allow recording the current streams.
|
||||
# pipe.transformer.enable_group_offload(onload_device=onload_device, offload_device=offload_device, offload_type="leaf_level", use_stream=True, record_stream=True)
|
||||
|
||||
# For any other model implementations, the apply_group_offloading function can be used
|
||||
apply_group_offloading(pipe.text_encoder, onload_device=onload_device, offload_type="block_level", num_blocks_per_group=2)
|
||||
apply_group_offloading(pipe.vae, onload_device=onload_device, offload_type="leaf_level")
|
||||
|
||||
prompt = (
|
||||
"A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. "
|
||||
"The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other "
|
||||
"pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, "
|
||||
"casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. "
|
||||
"The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical "
|
||||
"atmosphere of this unique musical performance."
|
||||
)
|
||||
video = pipe(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
|
||||
# This utilized about 14.79 GB. It can be further reduced by using tiling and using leaf_level offloading throughout the pipeline.
|
||||
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
|
||||
export_to_video(video, "output.mp4", fps=8)
|
||||
```
|
||||
|
||||
Group offloading (for CUDA devices with support for asynchronous data transfer streams) overlaps data transfer and computation to reduce the overall execution time compared to sequential offloading. This is enabled using layer prefetching with CUDA streams. The next layer to be executed is loaded onto the accelerator device while the current layer is being executed - this increases the memory requirements slightly. Group offloading also supports leaf-level offloading (equivalent to sequential CPU offloading) but can be made much faster when using streams.
|
||||
|
||||
<Tip>
|
||||
|
||||
- Group offloading may not work with all models out-of-the-box. If the forward implementations of the model contain weight-dependent device-casting of inputs, it may clash with the offloading mechanism's handling of device-casting.
|
||||
- The `offload_type` parameter can be set to either `block_level` or `leaf_level`. `block_level` offloads groups of `torch::nn::ModuleList` or `torch::nn:Sequential` modules based on a configurable attribute `num_blocks_per_group`. For example, if you set `num_blocks_per_group=2` on a standard transformer model containing 40 layers, it will onload/offload 2 layers at a time for a total of 20 onload/offloads. This drastically reduces the VRAM requirements. `leaf_level` offloads individual layers at the lowest level, which is equivalent to sequential offloading. However, unlike sequential offloading, group offloading can be made much faster when using streams, with minimal compromise to end-to-end generation time.
|
||||
- The `use_stream` parameter can be used with CUDA devices to enable prefetching layers for onload. It defaults to `False`. Layer prefetching allows overlapping computation and data transfer of model weights, which drastically reduces the overall execution time compared to other offloading methods. However, it can increase the CPU RAM usage significantly. Ensure that available CPU RAM that is at least twice the size of the model when setting `use_stream=True`. You can find more information about CUDA streams [here](https://pytorch.org/docs/stable/generated/torch.cuda.Stream.html)
|
||||
- If specifying `use_stream=True` on VAEs with tiling enabled, make sure to do a dummy forward pass (possibly with dummy inputs) before the actual inference to avoid device-mismatch errors. This may not work on all implementations. Please open an issue if you encounter any problems.
|
||||
- The parameter `low_cpu_mem_usage` can be set to `True` to reduce CPU memory usage when using streams for group offloading. This is useful when the CPU memory is the bottleneck, but it may counteract the benefits of using streams and increase the overall execution time. The CPU memory savings come from creating pinned-tensors on-the-fly instead of pre-pinning them. This parameter is better suited for using `leaf_level` offloading.
|
||||
- When using `use_stream=True`, users can additionally specify `record_stream=True` to get better speedups at the expense of slightly increased memory usage. Refer to the [official PyTorch docs](https://pytorch.org/docs/stable/generated/torch.Tensor.record_stream.html) to know more about this.
|
||||
|
||||
For more information about available parameters and an explanation of how group offloading works, refer to [`~hooks.group_offloading.apply_group_offloading`].
|
||||
|
||||
</Tip>
|
||||
|
||||
## FP8 layerwise weight-casting
|
||||
|
||||
PyTorch supports `torch.float8_e4m3fn` and `torch.float8_e5m2` as weight storage dtypes, but they can't be used for computation in many different tensor operations due to unimplemented kernel support. However, you can use these dtypes to store model weights in fp8 precision and upcast them on-the-fly when the layers are used in the forward pass. This is known as layerwise weight-casting.
|
||||
|
||||
Typically, inference on most models is done with `torch.float16` or `torch.bfloat16` weight/computation precision. Layerwise weight-casting cuts down the memory footprint of the model weights by approximately half.
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import CogVideoXPipeline, CogVideoXTransformer3DModel
|
||||
from diffusers.utils import export_to_video
|
||||
|
||||
model_id = "THUDM/CogVideoX-5b"
|
||||
|
||||
# Load the model in bfloat16 and enable layerwise casting
|
||||
transformer = CogVideoXTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
|
||||
transformer.enable_layerwise_casting(storage_dtype=torch.float8_e4m3fn, compute_dtype=torch.bfloat16)
|
||||
|
||||
# Load the pipeline
|
||||
pipe = CogVideoXPipeline.from_pretrained(model_id, transformer=transformer, torch_dtype=torch.bfloat16)
|
||||
pipe.to("cuda")
|
||||
|
||||
prompt = (
|
||||
"A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. "
|
||||
"The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other "
|
||||
"pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, "
|
||||
"casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. "
|
||||
"The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical "
|
||||
"atmosphere of this unique musical performance."
|
||||
)
|
||||
video = pipe(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
|
||||
export_to_video(video, "output.mp4", fps=8)
|
||||
```
|
||||
|
||||
In the above example, layerwise casting is enabled on the transformer component of the pipeline. By default, certain layers are skipped from the FP8 weight casting because it can lead to significant degradation of generation quality. The normalization and modulation related weight parameters are also skipped by default.
|
||||
|
||||
However, you gain more control and flexibility by directly utilizing the [`~hooks.layerwise_casting.apply_layerwise_casting`] function instead of [`~ModelMixin.enable_layerwise_casting`].
|
||||
|
||||
<Tip>
|
||||
|
||||
- Layerwise casting may not work with all models out-of-the-box. Sometimes, the forward implementations of the model might contain internal typecasting of weight values. Such implementations are not supported due to the currently simplistic implementation of layerwise casting, which assumes that the forward pass is independent of the weight precision and that the input dtypes are always in `compute_dtype`. An example of an incompatible implementation can be found [here](https://github.com/huggingface/transformers/blob/7f5077e53682ca855afc826162b204ebf809f1f9/src/transformers/models/t5/modeling_t5.py#L294-L299).
|
||||
- Layerwise casting may fail on custom modeling implementations that make use of [PEFT](https://github.com/huggingface/peft) layers. Some minimal checks to handle this case is implemented but is not extensively tested or guaranteed to work in all cases.
|
||||
- It can be also be applied partially to specific layers of a model. Partially applying layerwise casting can either be done manually by calling the `apply_layerwise_casting` function on specific internal modules, or by specifying the `skip_modules_pattern` and `skip_modules_classes` parameters for a root module. These parameters are particularly useful for layers such as normalization and modulation.
|
||||
|
||||
</Tip>
|
||||
|
||||
## Channels-last memory format
|
||||
|
||||
The channels-last memory format is an alternative way of ordering NCHW tensors in memory to preserve dimension ordering. Channels-last tensors are ordered in such a way that the channels become the densest dimension (storing images pixel-per-pixel). Since not all operators currently support the channels-last format, it may result in worst performance but you should still try and see if it works for your model.
|
||||
|
||||
@@ -12,6 +12,9 @@ specific language governing permissions and limitations under the License.
|
||||
|
||||
# Metal Performance Shaders (MPS)
|
||||
|
||||
> [!TIP]
|
||||
> Pipelines with a <img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22"> badge indicate a model can take advantage of the MPS backend on Apple silicon devices for faster inference. Feel free to open a [Pull Request](https://github.com/huggingface/diffusers/compare) to add this badge to pipelines that are missing it.
|
||||
|
||||
🤗 Diffusers is compatible with Apple silicon (M1/M2 chips) using the PyTorch [`mps`](https://pytorch.org/docs/stable/notes/mps.html) device, which uses the Metal framework to leverage the GPU on MacOS devices. You'll need to have:
|
||||
|
||||
- macOS computer with Apple silicon (M1/M2) hardware
|
||||
@@ -37,7 +40,7 @@ image
|
||||
|
||||
<Tip warning={true}>
|
||||
|
||||
Generating multiple prompts in a batch can [crash](https://github.com/huggingface/diffusers/issues/363) or fail to work reliably. We believe this is related to the [`mps`](https://github.com/pytorch/pytorch/issues/84039) backend in PyTorch. While this is being investigated, you should iterate instead of batching.
|
||||
The PyTorch [mps](https://pytorch.org/docs/stable/notes/mps.html) backend does not support NDArray sizes greater than `2**32`. Please open an [Issue](https://github.com/huggingface/diffusers/issues/new/choose) if you encounter this problem so we can investigate.
|
||||
|
||||
</Tip>
|
||||
|
||||
@@ -59,6 +62,10 @@ If you're using **PyTorch 1.13**, you need to "prime" the pipeline with an addit
|
||||
|
||||
## Troubleshoot
|
||||
|
||||
This section lists some common issues with using the `mps` backend and how to solve them.
|
||||
|
||||
### Attention slicing
|
||||
|
||||
M1/M2 performance is very sensitive to memory pressure. When this occurs, the system automatically swaps if it needs to which significantly degrades performance.
|
||||
|
||||
To prevent this from happening, we recommend *attention slicing* to reduce memory pressure during inference and prevent swapping. This is especially relevant if your computer has less than 64GB of system RAM, or if you generate images at non-standard resolutions larger than 512×512 pixels. Call the [`~DiffusionPipeline.enable_attention_slicing`] function on your pipeline:
|
||||
@@ -72,3 +79,7 @@ pipeline.enable_attention_slicing()
|
||||
```
|
||||
|
||||
Attention slicing performs the costly attention operation in multiple steps instead of all at once. It usually improves performance by ~20% in computers without universal memory, but we've observed *better performance* in most Apple silicon computers unless you have 64GB of RAM or more.
|
||||
|
||||
### Batch inference
|
||||
|
||||
Generating multiple prompts in a batch can crash or fail to work reliably. If this is the case, try iterating instead of batching.
|
||||
@@ -29,7 +29,7 @@ However, it is hard to decide when to reuse the cache to ensure quality generate
|
||||
This achieves a 2x speedup on FLUX.1-dev and HunyuanVideo inference with very good quality.
|
||||
|
||||
<figure>
|
||||
<img src="https://huggingface.co/datasets/chengzeyi/documentation-images/resolve/main/diffusers/para-attn/ada-cache.png" alt="Cache in Diffusion Transformer" />
|
||||
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/ada-cache.png" alt="Cache in Diffusion Transformer" />
|
||||
<figcaption>How AdaCache works, First Block Cache is a variant of it</figcaption>
|
||||
</figure>
|
||||
|
||||
|
||||
@@ -49,7 +49,7 @@ For Ada and higher-series GPUs. we recommend changing `torch_dtype` to `torch.bf
|
||||
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
|
||||
from transformers import BitsAndBytesConfig as TransformersBitsAndBytesConfig
|
||||
|
||||
from diffusers import FluxTransformer2DModel
|
||||
from diffusers import AutoModel
|
||||
from transformers import T5EncoderModel
|
||||
|
||||
quant_config = TransformersBitsAndBytesConfig(load_in_8bit=True,)
|
||||
@@ -63,7 +63,7 @@ text_encoder_2_8bit = T5EncoderModel.from_pretrained(
|
||||
|
||||
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True,)
|
||||
|
||||
transformer_8bit = FluxTransformer2DModel.from_pretrained(
|
||||
transformer_8bit = AutoModel.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-dev",
|
||||
subfolder="transformer",
|
||||
quantization_config=quant_config,
|
||||
@@ -74,7 +74,7 @@ transformer_8bit = FluxTransformer2DModel.from_pretrained(
|
||||
By default, all the other modules such as `torch.nn.LayerNorm` are converted to `torch.float16`. You can change the data type of these modules with the `torch_dtype` parameter.
|
||||
|
||||
```diff
|
||||
transformer_8bit = FluxTransformer2DModel.from_pretrained(
|
||||
transformer_8bit = AutoModel.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-dev",
|
||||
subfolder="transformer",
|
||||
quantization_config=quant_config,
|
||||
@@ -133,7 +133,7 @@ For Ada and higher-series GPUs. we recommend changing `torch_dtype` to `torch.bf
|
||||
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
|
||||
from transformers import BitsAndBytesConfig as TransformersBitsAndBytesConfig
|
||||
|
||||
from diffusers import FluxTransformer2DModel
|
||||
from diffusers import AutoModel
|
||||
from transformers import T5EncoderModel
|
||||
|
||||
quant_config = TransformersBitsAndBytesConfig(load_in_4bit=True,)
|
||||
@@ -147,7 +147,7 @@ text_encoder_2_4bit = T5EncoderModel.from_pretrained(
|
||||
|
||||
quant_config = DiffusersBitsAndBytesConfig(load_in_4bit=True,)
|
||||
|
||||
transformer_4bit = FluxTransformer2DModel.from_pretrained(
|
||||
transformer_4bit = AutoModel.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-dev",
|
||||
subfolder="transformer",
|
||||
quantization_config=quant_config,
|
||||
@@ -158,7 +158,7 @@ transformer_4bit = FluxTransformer2DModel.from_pretrained(
|
||||
By default, all the other modules such as `torch.nn.LayerNorm` are converted to `torch.float16`. You can change the data type of these modules with the `torch_dtype` parameter.
|
||||
|
||||
```diff
|
||||
transformer_4bit = FluxTransformer2DModel.from_pretrained(
|
||||
transformer_4bit = AutoModel.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-dev",
|
||||
subfolder="transformer",
|
||||
quantization_config=quant_config,
|
||||
@@ -217,11 +217,11 @@ print(model.get_memory_footprint())
|
||||
Quantized models can be loaded from the [`~ModelMixin.from_pretrained`] method without needing to specify the `quantization_config` parameters:
|
||||
|
||||
```py
|
||||
from diffusers import FluxTransformer2DModel, BitsAndBytesConfig
|
||||
from diffusers import AutoModel, BitsAndBytesConfig
|
||||
|
||||
quantization_config = BitsAndBytesConfig(load_in_4bit=True)
|
||||
|
||||
model_4bit = FluxTransformer2DModel.from_pretrained(
|
||||
model_4bit = AutoModel.from_pretrained(
|
||||
"hf-internal-testing/flux.1-dev-nf4-pkg", subfolder="transformer"
|
||||
)
|
||||
```
|
||||
@@ -243,13 +243,13 @@ An "outlier" is a hidden state value greater than a certain threshold, and these
|
||||
To find the best threshold for your model, we recommend experimenting with the `llm_int8_threshold` parameter in [`BitsAndBytesConfig`]:
|
||||
|
||||
```py
|
||||
from diffusers import FluxTransformer2DModel, BitsAndBytesConfig
|
||||
from diffusers import AutoModel, BitsAndBytesConfig
|
||||
|
||||
quantization_config = BitsAndBytesConfig(
|
||||
load_in_8bit=True, llm_int8_threshold=10,
|
||||
)
|
||||
|
||||
model_8bit = FluxTransformer2DModel.from_pretrained(
|
||||
model_8bit = AutoModel.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-dev",
|
||||
subfolder="transformer",
|
||||
quantization_config=quantization_config,
|
||||
@@ -305,7 +305,7 @@ NF4 is a 4-bit data type from the [QLoRA](https://hf.co/papers/2305.14314) paper
|
||||
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
|
||||
from transformers import BitsAndBytesConfig as TransformersBitsAndBytesConfig
|
||||
|
||||
from diffusers import FluxTransformer2DModel
|
||||
from diffusers import AutoModel
|
||||
from transformers import T5EncoderModel
|
||||
|
||||
quant_config = TransformersBitsAndBytesConfig(
|
||||
@@ -325,7 +325,7 @@ quant_config = DiffusersBitsAndBytesConfig(
|
||||
bnb_4bit_quant_type="nf4",
|
||||
)
|
||||
|
||||
transformer_4bit = FluxTransformer2DModel.from_pretrained(
|
||||
transformer_4bit = AutoModel.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-dev",
|
||||
subfolder="transformer",
|
||||
quantization_config=quant_config,
|
||||
@@ -343,7 +343,7 @@ Nested quantization is a technique that can save additional memory at no additio
|
||||
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
|
||||
from transformers import BitsAndBytesConfig as TransformersBitsAndBytesConfig
|
||||
|
||||
from diffusers import FluxTransformer2DModel
|
||||
from diffusers import AutoModel
|
||||
from transformers import T5EncoderModel
|
||||
|
||||
quant_config = TransformersBitsAndBytesConfig(
|
||||
@@ -363,7 +363,7 @@ quant_config = DiffusersBitsAndBytesConfig(
|
||||
bnb_4bit_use_double_quant=True,
|
||||
)
|
||||
|
||||
transformer_4bit = FluxTransformer2DModel.from_pretrained(
|
||||
transformer_4bit = AutoModel.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-dev",
|
||||
subfolder="transformer",
|
||||
quantization_config=quant_config,
|
||||
@@ -379,7 +379,7 @@ Once quantized, you can dequantize a model to its original precision, but this m
|
||||
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
|
||||
from transformers import BitsAndBytesConfig as TransformersBitsAndBytesConfig
|
||||
|
||||
from diffusers import FluxTransformer2DModel
|
||||
from diffusers import AutoModel
|
||||
from transformers import T5EncoderModel
|
||||
|
||||
quant_config = TransformersBitsAndBytesConfig(
|
||||
@@ -399,7 +399,7 @@ quant_config = DiffusersBitsAndBytesConfig(
|
||||
bnb_4bit_use_double_quant=True,
|
||||
)
|
||||
|
||||
transformer_4bit = FluxTransformer2DModel.from_pretrained(
|
||||
transformer_4bit = AutoModel.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-dev",
|
||||
subfolder="transformer",
|
||||
quantization_config=quant_config,
|
||||
|
||||
@@ -36,5 +36,6 @@ Diffusers currently supports the following quantization methods.
|
||||
- [BitsandBytes](./bitsandbytes)
|
||||
- [TorchAO](./torchao)
|
||||
- [GGUF](./gguf)
|
||||
- [Quanto](./quanto.md)
|
||||
|
||||
[This resource](https://huggingface.co/docs/transformers/main/en/quantization/overview#when-to-use-what) provides a good overview of the pros and cons of different quantization techniques.
|
||||
|
||||
148
docs/source/en/quantization/quanto.md
Normal file
148
docs/source/en/quantization/quanto.md
Normal file
@@ -0,0 +1,148 @@
|
||||
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
|
||||
-->
|
||||
|
||||
# Quanto
|
||||
|
||||
[Quanto](https://github.com/huggingface/optimum-quanto) is a PyTorch quantization backend for [Optimum](https://huggingface.co/docs/optimum/en/index). It has been designed with versatility and simplicity in mind:
|
||||
|
||||
- All features are available in eager mode (works with non-traceable models)
|
||||
- Supports quantization aware training
|
||||
- Quantized models are compatible with `torch.compile`
|
||||
- Quantized models are Device agnostic (e.g CUDA,XPU,MPS,CPU)
|
||||
|
||||
In order to use the Quanto backend, you will first need to install `optimum-quanto>=0.2.6` and `accelerate`
|
||||
|
||||
```shell
|
||||
pip install optimum-quanto accelerate
|
||||
```
|
||||
|
||||
Now you can quantize a model by passing the `QuantoConfig` object to the `from_pretrained()` method. Although the Quanto library does allow quantizing `nn.Conv2d` and `nn.LayerNorm` modules, currently, Diffusers only supports quantizing the weights in the `nn.Linear` layers of a model. The following snippet demonstrates how to apply `float8` quantization with Quanto.
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import FluxTransformer2DModel, QuantoConfig
|
||||
|
||||
model_id = "black-forest-labs/FLUX.1-dev"
|
||||
quantization_config = QuantoConfig(weights_dtype="float8")
|
||||
transformer = FluxTransformer2DModel.from_pretrained(
|
||||
model_id,
|
||||
subfolder="transformer",
|
||||
quantization_config=quantization_config,
|
||||
torch_dtype=torch.bfloat16,
|
||||
)
|
||||
|
||||
pipe = FluxPipeline.from_pretrained(model_id, transformer=transformer, torch_dtype=torch_dtype)
|
||||
pipe.to("cuda")
|
||||
|
||||
prompt = "A cat holding a sign that says hello world"
|
||||
image = pipe(
|
||||
prompt, num_inference_steps=50, guidance_scale=4.5, max_sequence_length=512
|
||||
).images[0]
|
||||
image.save("output.png")
|
||||
```
|
||||
|
||||
## Skipping Quantization on specific modules
|
||||
|
||||
It is possible to skip applying quantization on certain modules using the `modules_to_not_convert` argument in the `QuantoConfig`. Please ensure that the modules passed in to this argument match the keys of the modules in the `state_dict`
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import FluxTransformer2DModel, QuantoConfig
|
||||
|
||||
model_id = "black-forest-labs/FLUX.1-dev"
|
||||
quantization_config = QuantoConfig(weights_dtype="float8", modules_to_not_convert=["proj_out"])
|
||||
transformer = FluxTransformer2DModel.from_pretrained(
|
||||
model_id,
|
||||
subfolder="transformer",
|
||||
quantization_config=quantization_config,
|
||||
torch_dtype=torch.bfloat16,
|
||||
)
|
||||
```
|
||||
|
||||
## Using `from_single_file` with the Quanto Backend
|
||||
|
||||
`QuantoConfig` is compatible with `~FromOriginalModelMixin.from_single_file`.
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import FluxTransformer2DModel, QuantoConfig
|
||||
|
||||
ckpt_path = "https://huggingface.co/black-forest-labs/FLUX.1-dev/blob/main/flux1-dev.safetensors"
|
||||
quantization_config = QuantoConfig(weights_dtype="float8")
|
||||
transformer = FluxTransformer2DModel.from_single_file(ckpt_path, quantization_config=quantization_config, torch_dtype=torch.bfloat16)
|
||||
```
|
||||
|
||||
## Saving Quantized models
|
||||
|
||||
Diffusers supports serializing Quanto models using the `~ModelMixin.save_pretrained` method.
|
||||
|
||||
The serialization and loading requirements are different for models quantized directly with the Quanto library and models quantized
|
||||
with Diffusers using Quanto as the backend. It is currently not possible to load models quantized directly with Quanto into Diffusers using `~ModelMixin.from_pretrained`
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import FluxTransformer2DModel, QuantoConfig
|
||||
|
||||
model_id = "black-forest-labs/FLUX.1-dev"
|
||||
quantization_config = QuantoConfig(weights_dtype="float8")
|
||||
transformer = FluxTransformer2DModel.from_pretrained(
|
||||
model_id,
|
||||
subfolder="transformer",
|
||||
quantization_config=quantization_config,
|
||||
torch_dtype=torch.bfloat16,
|
||||
)
|
||||
# save quantized model to reuse
|
||||
transformer.save_pretrained("<your quantized model save path>")
|
||||
|
||||
# you can reload your quantized model with
|
||||
model = FluxTransformer2DModel.from_pretrained("<your quantized model save path>")
|
||||
```
|
||||
|
||||
## Using `torch.compile` with Quanto
|
||||
|
||||
Currently the Quanto backend supports `torch.compile` for the following quantization types:
|
||||
|
||||
- `int8` weights
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import FluxPipeline, FluxTransformer2DModel, QuantoConfig
|
||||
|
||||
model_id = "black-forest-labs/FLUX.1-dev"
|
||||
quantization_config = QuantoConfig(weights_dtype="int8")
|
||||
transformer = FluxTransformer2DModel.from_pretrained(
|
||||
model_id,
|
||||
subfolder="transformer",
|
||||
quantization_config=quantization_config,
|
||||
torch_dtype=torch.bfloat16,
|
||||
)
|
||||
transformer = torch.compile(transformer, mode="max-autotune", fullgraph=True)
|
||||
|
||||
pipe = FluxPipeline.from_pretrained(
|
||||
model_id, transformer=transformer, torch_dtype=torch_dtype
|
||||
)
|
||||
pipe.to("cuda")
|
||||
images = pipe("A cat holding a sign that says hello").images[0]
|
||||
images.save("flux-quanto-compile.png")
|
||||
```
|
||||
|
||||
## Supported Quantization Types
|
||||
|
||||
### Weights
|
||||
|
||||
- float8
|
||||
- int8
|
||||
- int4
|
||||
- int2
|
||||
|
||||
|
||||
@@ -26,13 +26,13 @@ The example below only quantizes the weights to int8.
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import FluxPipeline, FluxTransformer2DModel, TorchAoConfig
|
||||
from diffusers import FluxPipeline, AutoModel, TorchAoConfig
|
||||
|
||||
model_id = "black-forest-labs/FLUX.1-dev"
|
||||
dtype = torch.bfloat16
|
||||
|
||||
quantization_config = TorchAoConfig("int8wo")
|
||||
transformer = FluxTransformer2DModel.from_pretrained(
|
||||
transformer = AutoModel.from_pretrained(
|
||||
model_id,
|
||||
subfolder="transformer",
|
||||
quantization_config=quantization_config,
|
||||
@@ -99,10 +99,10 @@ To serialize a quantized model in a given dtype, first load the model with the d
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import FluxTransformer2DModel, TorchAoConfig
|
||||
from diffusers import AutoModel, TorchAoConfig
|
||||
|
||||
quantization_config = TorchAoConfig("int8wo")
|
||||
transformer = FluxTransformer2DModel.from_pretrained(
|
||||
transformer = AutoModel.from_pretrained(
|
||||
"black-forest-labs/Flux.1-Dev",
|
||||
subfolder="transformer",
|
||||
quantization_config=quantization_config,
|
||||
@@ -115,9 +115,9 @@ To load a serialized quantized model, use the [`~ModelMixin.from_pretrained`] me
|
||||
|
||||
```python
|
||||
import torch
|
||||
from diffusers import FluxPipeline, FluxTransformer2DModel
|
||||
from diffusers import FluxPipeline, AutoModel
|
||||
|
||||
transformer = FluxTransformer2DModel.from_pretrained("/path/to/flux_int8wo", torch_dtype=torch.bfloat16, use_safetensors=False)
|
||||
transformer = AutoModel.from_pretrained("/path/to/flux_int8wo", torch_dtype=torch.bfloat16, use_safetensors=False)
|
||||
pipe = FluxPipeline.from_pretrained("black-forest-labs/Flux.1-Dev", transformer=transformer, torch_dtype=torch.bfloat16)
|
||||
pipe.to("cuda")
|
||||
|
||||
@@ -126,15 +126,15 @@ image = pipe(prompt, num_inference_steps=30, guidance_scale=7.0).images[0]
|
||||
image.save("output.png")
|
||||
```
|
||||
|
||||
Some quantization methods, such as `uint4wo`, cannot be loaded directly and may result in an `UnpicklingError` when trying to load the models, but work as expected when saving them. In order to work around this, one can load the state dict manually into the model. Note, however, that this requires using `weights_only=False` in `torch.load`, so it should be run only if the weights were obtained from a trustable source.
|
||||
If you are using `torch<=2.6.0`, some quantization methods, such as `uint4wo`, cannot be loaded directly and may result in an `UnpicklingError` when trying to load the models, but work as expected when saving them. In order to work around this, one can load the state dict manually into the model. Note, however, that this requires using `weights_only=False` in `torch.load`, so it should be run only if the weights were obtained from a trustable source.
|
||||
|
||||
```python
|
||||
import torch
|
||||
from accelerate import init_empty_weights
|
||||
from diffusers import FluxPipeline, FluxTransformer2DModel, TorchAoConfig
|
||||
from diffusers import FluxPipeline, AutoModel, TorchAoConfig
|
||||
|
||||
# Serialize the model
|
||||
transformer = FluxTransformer2DModel.from_pretrained(
|
||||
transformer = AutoModel.from_pretrained(
|
||||
"black-forest-labs/Flux.1-Dev",
|
||||
subfolder="transformer",
|
||||
quantization_config=TorchAoConfig("uint4wo"),
|
||||
@@ -146,10 +146,13 @@ transformer.save_pretrained("/path/to/flux_uint4wo", safe_serialization=False, m
|
||||
# Load the model
|
||||
state_dict = torch.load("/path/to/flux_uint4wo/diffusion_pytorch_model.bin", weights_only=False, map_location="cpu")
|
||||
with init_empty_weights():
|
||||
transformer = FluxTransformer2DModel.from_config("/path/to/flux_uint4wo/config.json")
|
||||
transformer = AutoModel.from_config("/path/to/flux_uint4wo/config.json")
|
||||
transformer.load_state_dict(state_dict, strict=True, assign=True)
|
||||
```
|
||||
|
||||
> [!TIP]
|
||||
> The [`AutoModel`] API is supported for PyTorch >= 2.6 as shown in the examples below.
|
||||
|
||||
## Resources
|
||||
|
||||
- [TorchAO Quantization API](https://github.com/pytorch/ao/blob/main/torchao/quantization/README.md)
|
||||
|
||||
@@ -163,6 +163,9 @@ Models are initiated with the [`~ModelMixin.from_pretrained`] method which also
|
||||
>>> model = UNet2DModel.from_pretrained(repo_id, use_safetensors=True)
|
||||
```
|
||||
|
||||
> [!TIP]
|
||||
> Use the [`AutoModel`] API to automatically select a model class if you're unsure of which one to use.
|
||||
|
||||
To access the model parameters, call `model.config`:
|
||||
|
||||
```py
|
||||
|
||||
@@ -31,10 +31,10 @@ To adapt your text-to-image model for inpainting, you'll need to change the numb
|
||||
Initialize a [`UNet2DConditionModel`] with the pretrained text-to-image model weights, and change `in_channels` to 9. Changing the number of `in_channels` means you need to set `ignore_mismatched_sizes=True` and `low_cpu_mem_usage=False` to avoid a size mismatch error because the shape is different now.
|
||||
|
||||
```py
|
||||
from diffusers import UNet2DConditionModel
|
||||
from diffusers import AutoModel
|
||||
|
||||
model_id = "stable-diffusion-v1-5/stable-diffusion-v1-5"
|
||||
unet = UNet2DConditionModel.from_pretrained(
|
||||
unet = AutoModel.from_pretrained(
|
||||
model_id,
|
||||
subfolder="unet",
|
||||
in_channels=9,
|
||||
|
||||
@@ -339,7 +339,10 @@ import torch
|
||||
from huggingface_hub.repocard import RepoCard
|
||||
from diffusers import DiffusionPipeline
|
||||
|
||||
pipeline = DiffusionPipeline.from_pretrained("sayakpaul/custom-diffusion-cat-wooden-pot", torch_dtype=torch.float16).to("cuda")
|
||||
pipeline = DiffusionPipeline.from_pretrained(
|
||||
"CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16,
|
||||
).to("cuda")
|
||||
model_id = "sayakpaul/custom-diffusion-cat-wooden-pot"
|
||||
pipeline.unet.load_attn_procs(model_id, weight_name="pytorch_custom_diffusion_weights.bin")
|
||||
pipeline.load_textual_inversion(model_id, weight_name="<new1>.bin")
|
||||
pipeline.load_textual_inversion(model_id, weight_name="<new2>.bin")
|
||||
|
||||
@@ -165,10 +165,10 @@ flush()
|
||||
Load the diffusion transformer next which has 12.5B parameters. This time, set `device_map="auto"` to automatically distribute the model across two 16GB GPUs. The `auto` strategy is backed by [Accelerate](https://hf.co/docs/accelerate/index) and available as a part of the [Big Model Inference](https://hf.co/docs/accelerate/concept_guides/big_model_inference) feature. It starts by distributing a model across the fastest device first (GPU) before moving to slower devices like the CPU and hard drive if needed. The trade-off of storing model parameters on slower devices is slower inference latency.
|
||||
|
||||
```py
|
||||
from diffusers import FluxTransformer2DModel
|
||||
from diffusers import AutoModel
|
||||
import torch
|
||||
|
||||
transformer = FluxTransformer2DModel.from_pretrained(
|
||||
transformer = AutoModel.from_pretrained(
|
||||
"black-forest-labs/FLUX.1-dev",
|
||||
subfolder="transformer",
|
||||
device_map="auto",
|
||||
|
||||
Some files were not shown because too many files have changed in this diff Show More
Reference in New Issue
Block a user