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15 Commits

Author SHA1 Message Date
Sayak Paul
ed91e8b3e6 Merge branch 'main' into fix/single-file-playground 2024-03-07 15:00:33 +05:30
DN6
a77e426877 update 2024-03-07 13:15:51 +05:30
Sayak Paul
c0a0ef5deb Merge branch 'main' into fix/single-file-playground 2024-03-07 11:49:17 +05:30
sayakpaul
9e35a12587 fix 2024-03-06 18:09:05 +05:30
sayakpaul
49b0b516ea fix: kwargs 2024-03-06 18:04:22 +05:30
sayakpaul
52ba8061d3 address rest of the comments. 2024-03-06 18:00:37 +05:30
sayakpaul
2be231cce5 address Dhruv's comment. 2024-03-06 17:53:19 +05:30
Sayak Paul
4b315f16a8 Merge branch 'main' into fix/single-file-playground 2024-03-06 16:36:59 +05:30
sayakpaul
29e6b873c4 delegate model_type inference to a function. 2024-03-06 16:35:37 +05:30
sayakpaul
6d3e82c9cd fix: things. 2024-03-06 14:54:34 +05:30
Sayak Paul
1f358e1331 Merge branch 'main' into fix/single-file-playground 2024-03-06 14:41:17 +05:30
sayakpaul
c1d0e091af apply Dhruv's comments but errors. 2024-03-06 14:40:50 +05:30
sayakpaul
9d90d60753 fix: edm key 2024-03-06 13:10:49 +05:30
sayakpaul
a4e00abb68 remove is_playground_model. 2024-03-06 13:09:02 +05:30
sayakpaul
ce4f4f4545 fix: support for loading playground v2.5 single file checkpoint. 2024-03-06 11:41:16 +05:30
817 changed files with 19694 additions and 87457 deletions

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@@ -31,6 +31,7 @@ jobs:
nvidia-smi
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install pandas peft
@@ -39,7 +40,7 @@ jobs:
python utils/print_env.py
- name: Diffusers Benchmarking
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_BOT_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.DIFFUSERS_BOT_TOKEN }}
BASE_PATH: benchmark_outputs
run: |
export TOTAL_GPU_MEMORY=$(python -c "import torch; print(torch.cuda.get_device_properties(0).total_memory / (1024**3))")

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@@ -1,57 +1,21 @@
name: Test, build, and push Docker images
name: Build Docker images (nightly)
on:
pull_request: # During PRs, we just check if the changes Dockerfiles can be successfully built
branches:
- main
paths:
- "docker/**"
workflow_dispatch:
schedule:
- cron: "0 0 * * *" # every day at midnight
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
group: docker-image-builds
cancel-in-progress: false
env:
REGISTRY: diffusers
CI_SLACK_CHANNEL: ${{ secrets.CI_DOCKER_CHANNEL }}
jobs:
test-build-docker-images:
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
if: github.event_name == 'pull_request'
steps:
- name: Set up Docker Buildx
uses: docker/setup-buildx-action@v1
- name: Check out code
uses: actions/checkout@v3
- name: Find Changed Dockerfiles
id: file_changes
uses: jitterbit/get-changed-files@v1
with:
format: 'space-delimited'
token: ${{ secrets.GITHUB_TOKEN }}
- name: Build Changed Docker Images
run: |
CHANGED_FILES="${{ steps.file_changes.outputs.all }}"
for FILE in $CHANGED_FILES; do
if [[ "$FILE" == docker/*Dockerfile ]]; then
DOCKER_PATH="${FILE%/Dockerfile}"
DOCKER_TAG=$(basename "$DOCKER_PATH")
echo "Building Docker image for $DOCKER_TAG"
docker build -t "$DOCKER_TAG" "$DOCKER_PATH"
fi
done
if: steps.file_changes.outputs.all != ''
build-and-push-docker-images:
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
if: github.event_name != 'pull_request'
build-docker-images:
runs-on: ubuntu-latest
permissions:
contents: read
@@ -69,18 +33,17 @@ jobs:
- diffusers-flax-tpu
- diffusers-onnxruntime-cpu
- diffusers-onnxruntime-cuda
- diffusers-doc-builder
steps:
- name: Checkout repository
uses: actions/checkout@v3
- name: Set up Docker Buildx
uses: docker/setup-buildx-action@v1
- name: Login to Docker Hub
uses: docker/login-action@v2
with:
username: ${{ env.REGISTRY }}
password: ${{ secrets.DOCKERHUB_TOKEN }}
- name: Build and push
uses: docker/build-push-action@v3
with:
@@ -91,11 +54,24 @@ jobs:
- name: Post to a Slack channel
id: slack
uses: huggingface/hf-workflows/.github/actions/post-slack@main
uses: slackapi/slack-github-action@6c661ce58804a1a20f6dc5fbee7f0381b469e001
with:
# Slack channel id, channel name, or user id to post message.
# See also: https://api.slack.com/methods/chat.postMessage#channels
slack_channel: ${{ env.CI_SLACK_CHANNEL }}
title: "🤗 Results of the ${{ matrix.image-name }} Docker Image build"
status: ${{ job.status }}
slack_token: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
channel-id: ${{ env.CI_SLACK_CHANNEL }}
# For posting a rich message using Block Kit
payload: |
{
"text": "${{ matrix.image-name }} Docker Image build result: ${{ job.status }}\n${{ github.event.head_commit.url }}",
"blocks": [
{
"type": "section",
"text": {
"type": "mrkdwn",
"text": "${{ matrix.image-name }} Docker Image build result: ${{ job.status }}\n${{ github.event.head_commit.url }}"
}
}
]
}
env:
SLACK_BOT_TOKEN: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}

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@@ -21,7 +21,7 @@ jobs:
package: diffusers
notebook_folder: diffusers_doc
languages: en ko zh ja pt
custom_container: diffusers/diffusers-doc-builder
secrets:
token: ${{ secrets.HUGGINGFACE_PUSH }}
hf_token: ${{ secrets.HF_DOC_BUILD_PUSH }}

View File

@@ -20,4 +20,3 @@ jobs:
install_libgl1: true
package: diffusers
languages: en ko zh ja pt
custom_container: diffusers/diffusers-doc-builder

View File

@@ -1,89 +0,0 @@
name: Mirror Community Pipeline
on:
# Push changes on the main branch
push:
branches:
- main
paths:
- 'examples/community/**.py'
# And on tag creation (e.g. `v0.28.1`)
tags:
- '*'
# Manual trigger with ref input
workflow_dispatch:
inputs:
ref:
description: "Either 'main' or a tag ref"
required: true
default: 'main'
jobs:
mirror_community_pipeline:
runs-on: ubuntu-latest
steps:
# Checkout to correct ref
# If workflow dispatch
# If ref is 'main', set:
# CHECKOUT_REF=refs/heads/main
# PATH_IN_REPO=main
# Else it must be a tag. Set:
# CHECKOUT_REF=refs/tags/{tag}
# PATH_IN_REPO={tag}
# If not workflow dispatch
# If ref is 'refs/heads/main' => set 'main'
# Else it must be a tag => set {tag}
- name: Set checkout_ref and path_in_repo
run: |
if [ "${{ github.event_name }}" == "workflow_dispatch" ]; then
if [ -z "${{ github.event.inputs.ref }}" ]; then
echo "Error: Missing ref input"
exit 1
elif [ "${{ github.event.inputs.ref }}" == "main" ]; then
echo "CHECKOUT_REF=refs/heads/main" >> $GITHUB_ENV
echo "PATH_IN_REPO=main" >> $GITHUB_ENV
else
echo "CHECKOUT_REF=refs/tags/${{ github.event.inputs.ref }}" >> $GITHUB_ENV
echo "PATH_IN_REPO=${{ github.event.inputs.ref }}" >> $GITHUB_ENV
fi
elif [ "${{ github.ref }}" == "refs/heads/main" ]; then
echo "CHECKOUT_REF=${{ github.ref }}" >> $GITHUB_ENV
echo "PATH_IN_REPO=main" >> $GITHUB_ENV
else
# e.g. refs/tags/v0.28.1 -> v0.28.1
echo "CHECKOUT_REF=${{ github.ref }}" >> $GITHUB_ENV
echo "PATH_IN_REPO=$(echo ${{ github.ref }} | sed 's/^refs\/tags\///')" >> $GITHUB_ENV
fi
- name: Print env vars
run: |
echo "CHECKOUT_REF: ${{ env.CHECKOUT_REF }}"
echo "PATH_IN_REPO: ${{ env.PATH_IN_REPO }}"
- uses: actions/checkout@v3
with:
ref: ${{ env.CHECKOUT_REF }}
# Setup + install dependencies
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.10"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install --upgrade huggingface_hub
# Check secret is set
- name: whoami
run: huggingface-cli whoami
env:
HF_TOKEN: ${{ secrets.HF_TOKEN_MIRROR_COMMUNITY_PIPELINES }}
# Push to HF! (under subfolder based on checkout ref)
# https://huggingface.co/datasets/diffusers/community-pipelines-mirror
- name: Mirror community pipeline to HF
run: huggingface-cli upload diffusers/community-pipelines-mirror ./examples/community ${PATH_IN_REPO} --repo-type dataset
env:
PATH_IN_REPO: ${{ env.PATH_IN_REPO }}
HF_TOKEN: ${{ secrets.HF_TOKEN_MIRROR_COMMUNITY_PIPELINES }}

View File

@@ -1,7 +1,6 @@
name: Nightly and release tests on main/release branch
name: Nightly tests on main
on:
workflow_dispatch:
schedule:
- cron: "0 0 * * *" # every day at midnight
@@ -13,348 +12,110 @@ env:
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
RUN_NIGHTLY: yes
PIPELINE_USAGE_CUTOFF: 5000
SLACK_API_TOKEN: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines Matrix
runs-on: diffusers/diffusers-pytorch-cpu
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
pip install -e .
pip install huggingface_hub
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
run: |
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
echo $matrix
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: test-pipelines.json
path: reports
run_nightly_tests_for_torch_pipelines:
name: Torch Pipelines CUDA Nightly Tests
needs: setup_torch_cuda_pipeline_matrix
run_nightly_tests:
strategy:
fail-fast: false
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on: [single-gpu, nvidia-gpu, t4, ci]
config:
- name: Nightly PyTorch CUDA tests on Ubuntu
framework: pytorch
runner: docker-gpu
image: diffusers/diffusers-pytorch-cuda
report: torch_cuda
- name: Nightly Flax TPU tests on Ubuntu
framework: flax
runner: docker-tpu
image: diffusers/diffusers-flax-tpu
report: flax_tpu
- name: Nightly ONNXRuntime CUDA tests on Ubuntu
framework: onnxruntime
runner: docker-gpu
image: diffusers/diffusers-onnxruntime-cuda
report: onnx_cuda
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
image: ${{ matrix.config.image }}
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ ${{ matrix.config.runner == 'docker-tpu' && '--privileged' || '--gpus 0'}}
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: nvidia-smi
if: ${{ matrix.config.runner == 'docker-gpu' }}
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate
- name: Environment
run: |
python utils/print_env.py
- name: Nightly PyTorch CUDA checkpoint (pipelines) tests
- name: Run nightly PyTorch CUDA tests
if: ${{ matrix.config.framework == 'pytorch' }}
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
--report-log=tests_pipeline_${{ matrix.module }}_cuda.log \
tests/pipelines/${{ matrix.module }}
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Run nightly Flax TPU tests
if: ${{ matrix.config.framework == 'flax' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 0 \
-s -v -k "Flax" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Run nightly ONNXRuntime CUDA tests
if: ${{ matrix.config.framework == 'onnxruntime' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: pipeline_${{ matrix.module }}_test_reports
name: ${{ matrix.config.report }}_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_tests_for_other_torch_modules:
name: Torch Non-Pipelines CUDA Nightly Tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
strategy:
matrix:
module: [models, schedulers, others, examples]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly PyTorch CUDA tests for non-pipeline modules
if: ${{ matrix.module != 'examples'}}
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_${{ matrix.module }}_cuda \
--report-log=tests_torch_${{ matrix.module }}_cuda.log \
tests/${{ matrix.module }}
- name: Run nightly example tests with Torch
if: ${{ matrix.module == 'examples' }}
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v --make-reports=examples_torch_cuda \
--report-log=examples_torch_cuda.log \
examples/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_torch_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_${{ matrix.module }}_cuda_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_lora_nightly_tests:
name: Nightly LoRA Tests with PEFT and TORCH
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly LoRA tests with PEFT and Torch
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_lora_cuda \
--report-log=tests_torch_lora_cuda.log \
tests/lora
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_lora_cuda_stats.txt
cat reports/tests_torch_lora_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_lora_cuda_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_flax_tpu_tests:
name: Nightly Flax TPU Tests
runs-on: docker-tpu
if: github.event_name == 'schedule'
container:
image: diffusers/diffusers-flax-tpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --privileged
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly Flax TPU tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
--make-reports=tests_flax_tpu \
--report-log=tests_flax_tpu.log \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_flax_tpu_stats.txt
cat reports/tests_flax_tpu_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: flax_tpu_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_onnx_tests:
name: Nightly ONNXRuntime CUDA tests on Ubuntu
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-onnxruntime-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly ONNXRuntime CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_onnx_cuda \
--report-log=tests_onnx_cuda.log \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_onnx_cuda_stats.txt
cat reports/tests_onnx_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: ${{ matrix.config.report }}_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_tests_apple_m1:
name: Nightly PyTorch MPS tests on MacOS
runs-on: [ self-hosted, apple-m1 ]
if: github.event_name == 'schedule'
steps:
- name: Checkout diffusers
@@ -379,7 +140,6 @@ jobs:
${CONDA_RUN} python -m uv pip install -e [quality,test]
${CONDA_RUN} python -m uv pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
${CONDA_RUN} python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate
${CONDA_RUN} python -m uv pip install pytest-reportlog
- name: Environment
shell: arch -arch arm64 bash {0}
@@ -390,11 +150,9 @@ jobs:
shell: arch -arch arm64 bash {0}
env:
HF_HOME: /System/Volumes/Data/mnt/cache
HF_TOKEN: ${{ secrets.HF_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
--report-log=tests_torch_mps.log \
tests/
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps tests/
- name: Failure short reports
if: ${{ failure() }}
@@ -406,9 +164,3 @@ jobs:
with:
name: torch_mps_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY

View File

@@ -1,23 +0,0 @@
name: Notify Slack about a release
on:
workflow_dispatch:
release:
types: [published]
jobs:
build:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Setup Python
uses: actions/setup-python@v4
with:
python-version: '3.8'
- name: Notify Slack about the release
env:
SLACK_WEBHOOK_URL: ${{ secrets.SLACK_WEBHOOK_URL }}
run: pip install requests && python utils/notify_slack_about_release.py

View File

@@ -15,7 +15,7 @@ concurrency:
jobs:
setup_pr_tests:
name: Setup PR Tests
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on: docker-cpu
container:
image: diffusers/diffusers-pytorch-cpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
@@ -32,6 +32,7 @@ jobs:
fetch-depth: 0
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
- name: Environment
@@ -73,7 +74,7 @@ jobs:
max-parallel: 2
matrix:
modules: ${{ fromJson(needs.setup_pr_tests.outputs.matrix) }}
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on: docker-cpu
container:
image: diffusers/diffusers-pytorch-cpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
@@ -88,6 +89,7 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pip install -e [quality,test]
python -m pip install accelerate
@@ -123,7 +125,7 @@ jobs:
config:
- name: Hub tests for models, schedulers, and pipelines
framework: hub_tests_pytorch
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_hub
@@ -145,6 +147,7 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pip install -e [quality,test]

View File

@@ -32,11 +32,9 @@ jobs:
python -m pip install --upgrade pip
pip install .[quality]
- name: Check quality
run: make quality
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Quality check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make style && make quality'" >> $GITHUB_STEP_SUMMARY
ruff check examples tests src utils scripts
ruff format examples tests src utils scripts --check
check_repository_consistency:
needs: check_code_quality
@@ -51,15 +49,11 @@ jobs:
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check repo consistency
- name: Check quality
run: |
python utils/check_copies.py
python utils/check_dummies.py
make deps_table_check_updated
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Repo consistency check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make fix-copies'" >> $GITHUB_STEP_SUMMARY
run_fast_tests:
needs: [check_code_quality, check_repository_consistency]
@@ -71,7 +65,7 @@ jobs:
name: LoRA - ${{ matrix.lib-versions }}
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on: docker-cpu
container:
image: diffusers/diffusers-pytorch-cpu
@@ -89,10 +83,11 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
if [ "${{ matrix.lib-versions }}" == "main" ]; then
python -m pip install -U peft@git+https://github.com/huggingface/peft.git
python -m uv pip install -U peft@git+https://github.com/huggingface/peft.git
python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git
python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
else
@@ -107,25 +102,7 @@ jobs:
- name: Run fast PyTorch LoRA CPU tests with PEFT backend
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v \
--make-reports=tests_${{ matrix.config.report }} \
tests/lora/
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
-s -v \
--make-reports=tests_models_lora_${{ matrix.config.report }} \
tests/models/ -k "lora"
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_${{ matrix.config.report }}_failures_short.txt
cat reports/tests_models_lora_${{ matrix.config.report }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: pr_${{ matrix.config.report }}_test_reports
path: reports
tests/lora/test_lora_layers_peft.py

View File

@@ -40,11 +40,9 @@ jobs:
python -m pip install --upgrade pip
pip install .[quality]
- name: Check quality
run: make quality
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Quality check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make style && make quality'" >> $GITHUB_STEP_SUMMARY
ruff check examples tests src utils scripts
ruff format examples tests src utils scripts --check
check_repository_consistency:
needs: check_code_quality
@@ -59,15 +57,11 @@ jobs:
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check repo consistency
- name: Check quality
run: |
python utils/check_copies.py
python utils/check_dummies.py
make deps_table_check_updated
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Repo consistency check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make fix-copies'" >> $GITHUB_STEP_SUMMARY
run_fast_tests:
needs: [check_code_quality, check_repository_consistency]
@@ -77,22 +71,22 @@ jobs:
config:
- name: Fast PyTorch Pipeline CPU tests
framework: pytorch_pipelines
runner: [ self-hosted, intel-cpu, 32-cpu, 256-ram, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_pipelines
- name: Fast PyTorch Models & Schedulers CPU tests
framework: pytorch_models
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_models_schedulers
- name: Fast Flax CPU tests
framework: flax
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-flax-cpu
report: flax_cpu
- name: PyTorch Example CPU tests
framework: pytorch_examples
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_example_cpu
@@ -116,6 +110,7 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate
@@ -129,7 +124,7 @@ jobs:
if: ${{ matrix.config.framework == 'pytorch_pipelines' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 8 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/pipelines
@@ -138,7 +133,7 @@ jobs:
if: ${{ matrix.config.framework == 'pytorch_models' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx and not Dependency" \
--make-reports=tests_${{ matrix.config.report }} \
tests/models tests/schedulers tests/others
@@ -147,7 +142,7 @@ jobs:
if: ${{ matrix.config.framework == 'flax' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Flax" \
--make-reports=tests_${{ matrix.config.report }} \
tests
@@ -156,8 +151,8 @@ jobs:
if: ${{ matrix.config.framework == 'pytorch_examples' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install peft timm
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m uv pip install peft
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_${{ matrix.config.report }} \
examples
@@ -180,7 +175,7 @@ jobs:
config:
- name: Hub tests for models, schedulers, and pipelines
framework: hub_tests_pytorch
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_hub
@@ -204,6 +199,7 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]

View File

@@ -21,9 +21,10 @@ env:
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cpu
image: diffusers/diffusers-pytorch-cpu # this is a CPU image, but we need it to fetch the matrix
options: --shm-size "16gb" --ipc host
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
@@ -33,17 +34,22 @@ jobs:
fetch-depth: 2
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
run: |
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
echo $matrix
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
@@ -56,13 +62,12 @@ jobs:
needs: setup_torch_cuda_pipeline_matrix
strategy:
fail-fast: false
max-parallel: 8
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -73,6 +78,7 @@ jobs:
nvidia-smi
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
@@ -81,7 +87,7 @@ jobs:
python utils/print_env.py
- name: Slow PyTorch CUDA checkpoint tests on Ubuntu
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
@@ -94,6 +100,7 @@ jobs:
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
@@ -103,16 +110,16 @@ jobs:
torch_cuda_tests:
name: Torch CUDA Tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
strategy:
matrix:
module: [models, schedulers, lora, others, single_file]
module: [models, schedulers, lora, others]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -121,6 +128,7 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
@@ -131,7 +139,7 @@ jobs:
- name: Run slow PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
@@ -155,10 +163,10 @@ jobs:
peft_cuda_tests:
name: PEFT CUDA Tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
@@ -170,10 +178,11 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m pip install -U peft@git+https://github.com/huggingface/peft.git
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
- name: Environment
run: |
@@ -181,7 +190,7 @@ jobs:
- name: Run slow PEFT CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
@@ -189,17 +198,12 @@ jobs:
-s -v -k "not Flax and not Onnx and not PEFTLoRALoading" \
--make-reports=tests_peft_cuda \
tests/lora/
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "lora and not Flax and not Onnx and not PEFTLoRALoading" \
--make-reports=tests_peft_cuda_models_lora \
tests/models/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_peft_cuda_stats.txt
cat reports/tests_peft_cuda_failures_short.txt
cat reports/tests_peft_cuda_models_lora_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
@@ -213,7 +217,7 @@ jobs:
runs-on: docker-tpu
container:
image: diffusers/diffusers-flax-tpu
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/ --privileged
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --privileged
defaults:
run:
shell: bash
@@ -225,6 +229,7 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
@@ -235,7 +240,7 @@ jobs:
- name: Run slow Flax TPU tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
@@ -257,10 +262,10 @@ jobs:
onnx_cuda_tests:
name: ONNX CUDA Tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on: docker-gpu
container:
image: diffusers/diffusers-onnxruntime-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/ --gpus 0
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
@@ -272,6 +277,7 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
@@ -282,7 +288,7 @@ jobs:
- name: Run slow ONNXRuntime CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
@@ -305,11 +311,11 @@ jobs:
run_torch_compile_tests:
name: PyTorch Compile CUDA tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-compile-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
steps:
- name: Checkout diffusers
@@ -329,7 +335,7 @@ jobs:
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "compile" --make-reports=tests_torch_compile_cuda tests/
- name: Failure short reports
@@ -346,11 +352,11 @@ jobs:
run_xformers_tests:
name: PyTorch xformers CUDA tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-xformers-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
steps:
- name: Checkout diffusers
@@ -370,7 +376,7 @@ jobs:
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "xformers" --make-reports=tests_torch_xformers_cuda tests/
- name: Failure short reports
@@ -387,11 +393,11 @@ jobs:
run_examples_tests:
name: Examples PyTorch CUDA tests on Ubuntu
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
steps:
- name: Checkout diffusers
@@ -415,10 +421,9 @@ jobs:
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install timm
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=examples_torch_cuda examples/
- name: Failure short reports
@@ -432,4 +437,4 @@ jobs:
uses: actions/upload-artifact@v2
with:
name: examples_test_reports
path: reports
path: reports

View File

@@ -29,22 +29,22 @@ jobs:
config:
- name: Fast PyTorch CPU tests on Ubuntu
framework: pytorch
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu
- name: Fast Flax CPU tests on Ubuntu
framework: flax
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-flax-cpu
report: flax_cpu
- name: Fast ONNXRuntime CPU tests on Ubuntu
framework: onnxruntime
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-onnxruntime-cpu
report: onnx_cpu
- name: PyTorch Example CPU tests on Ubuntu
framework: pytorch_examples
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_example_cpu
@@ -68,6 +68,7 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
@@ -80,7 +81,7 @@ jobs:
if: ${{ matrix.config.framework == 'pytorch' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
@@ -89,7 +90,7 @@ jobs:
if: ${{ matrix.config.framework == 'flax' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Flax" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
@@ -98,7 +99,7 @@ jobs:
if: ${{ matrix.config.framework == 'onnxruntime' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
@@ -107,8 +108,8 @@ jobs:
if: ${{ matrix.config.framework == 'pytorch_examples' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install peft timm
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m uv pip install peft
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_${{ matrix.config.report }} \
examples

View File

@@ -23,7 +23,7 @@ concurrency:
jobs:
run_fast_tests_apple_m1:
name: Fast PyTorch MPS tests on MacOS
runs-on: macos-13-xlarge
runs-on: [ self-hosted, apple-m1 ]
steps:
- name: Checkout diffusers
@@ -59,7 +59,7 @@ jobs:
shell: arch -arch arm64 bash {0}
env:
HF_HOME: /System/Volumes/Data/mnt/cache
HF_TOKEN: ${{ secrets.HF_TOKEN }}
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
${CONDA_RUN} python -m pytest -n 0 -s -v --make-reports=tests_torch_mps tests/

View File

@@ -1,81 +0,0 @@
# Adapted from https://blog.deepjyoti30.dev/pypi-release-github-action
name: PyPI release
on:
workflow_dispatch:
push:
tags:
- "*"
jobs:
find-and-checkout-latest-branch:
runs-on: ubuntu-latest
outputs:
latest_branch: ${{ steps.set_latest_branch.outputs.latest_branch }}
steps:
- name: Checkout Repo
uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: '3.8'
- name: Fetch latest branch
id: fetch_latest_branch
run: |
pip install -U requests packaging
LATEST_BRANCH=$(python utils/fetch_latest_release_branch.py)
echo "Latest branch: $LATEST_BRANCH"
echo "latest_branch=$LATEST_BRANCH" >> $GITHUB_ENV
- name: Set latest branch output
id: set_latest_branch
run: echo "::set-output name=latest_branch::${{ env.latest_branch }}"
release:
needs: find-and-checkout-latest-branch
runs-on: ubuntu-latest
steps:
- name: Checkout Repo
uses: actions/checkout@v3
with:
ref: ${{ needs.find-and-checkout-latest-branch.outputs.latest_branch }}
- name: Setup Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install -U setuptools wheel twine
pip install -U torch --index-url https://download.pytorch.org/whl/cpu
pip install -U transformers
- name: Build the dist files
run: python setup.py bdist_wheel && python setup.py sdist
- name: Publish to the test PyPI
env:
TWINE_USERNAME: ${{ secrets.TEST_PYPI_USERNAME }}
TWINE_PASSWORD: ${{ secrets.TEST_PYPI_PASSWORD }}
run: twine upload dist/* -r pypitest --repository-url=https://test.pypi.org/legacy/
- name: Test installing diffusers and importing
run: |
pip install diffusers && pip uninstall diffusers -y
pip install -i https://testpypi.python.org/pypi diffusers
python -c "from diffusers import __version__; print(__version__)"
python -c "from diffusers import DiffusionPipeline; pipe = DiffusionPipeline.from_pretrained('fusing/unet-ldm-dummy-update'); pipe()"
python -c "from diffusers import DiffusionPipeline; pipe = DiffusionPipeline.from_pretrained('hf-internal-testing/tiny-stable-diffusion-pipe', safety_checker=None); pipe('ah suh du')"
python -c "from diffusers import *"
- name: Publish to PyPI
env:
TWINE_USERNAME: ${{ secrets.PYPI_USERNAME }}
TWINE_PASSWORD: ${{ secrets.PYPI_PASSWORD }}
run: twine upload dist/* -r pypi

View File

@@ -1,73 +0,0 @@
name: Check running SLOW tests from a PR (only GPU)
on:
workflow_dispatch:
inputs:
docker_image:
default: 'diffusers/diffusers-pytorch-cuda'
description: 'Name of the Docker image'
required: true
branch:
description: 'PR Branch to test on'
required: true
test:
description: 'Tests to run (e.g.: `tests/models`).'
required: true
env:
DIFFUSERS_IS_CI: yes
IS_GITHUB_CI: "1"
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
jobs:
run_tests:
name: "Run a test on our runner from a PR"
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: ${{ github.event.inputs.docker_image }}
options: --gpus 0 --privileged --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
steps:
- name: Validate test files input
id: validate_test_files
env:
PY_TEST: ${{ github.event.inputs.test }}
run: |
if [[ ! "$PY_TEST" =~ ^tests/ ]]; then
echo "Error: The input string must start with 'tests/'."
exit 1
fi
if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines) ]]; then
echo "Error: The input string must contain either 'models' or 'pipelines' after 'tests/'."
exit 1
fi
if [[ "$PY_TEST" == *";"* ]]; then
echo "Error: The input string must not contain ';'."
exit 1
fi
echo "$PY_TEST"
- name: Checkout PR branch
uses: actions/checkout@v4
with:
ref: ${{ github.event.inputs.branch }}
repository: ${{ github.event.pull_request.head.repo.full_name }}
- name: Install pytest
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install peft
- name: Run tests
env:
PY_TEST: ${{ github.event.inputs.test }}
run: |
pytest "$PY_TEST"

View File

@@ -1,46 +0,0 @@
name: SSH into runners
on:
workflow_dispatch:
inputs:
runner_type:
description: 'Type of runner to test (a10 or t4)'
required: true
docker_image:
description: 'Name of the Docker image'
required: true
env:
IS_GITHUB_CI: "1"
HF_HUB_READ_TOKEN: ${{ secrets.HF_HUB_READ_TOKEN }}
HF_HOME: /mnt/cache
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
RUN_SLOW: yes
jobs:
ssh_runner:
name: "SSH"
runs-on: [single-gpu, nvidia-gpu, "${{ github.event.inputs.runner_type }}", ci]
container:
image: ${{ github.event.inputs.docker_image }}
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0 --privileged
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Tailscale # In order to be able to SSH when a test fails
uses: huggingface/tailscale-action@main
with:
authkey: ${{ secrets.TAILSCALE_SSH_AUTHKEY }}
slackChannel: ${{ secrets.SLACK_CIFEEDBACK_CHANNEL }}
slackToken: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
waitForSSH: true

View File

@@ -1,15 +0,0 @@
on:
push:
name: Secret Leaks
jobs:
trufflehog:
runs-on: ubuntu-latest
steps:
- name: Checkout code
uses: actions/checkout@v4
with:
fetch-depth: 0
- name: Secret Scanning
uses: trufflesecurity/trufflehog@main

View File

@@ -1,30 +0,0 @@
name: Update Diffusers metadata
on:
workflow_dispatch:
push:
branches:
- main
- update_diffusers_metadata*
jobs:
update_metadata:
runs-on: ubuntu-22.04
defaults:
run:
shell: bash -l {0}
steps:
- uses: actions/checkout@v3
- name: Setup environment
run: |
pip install --upgrade pip
pip install datasets pandas
pip install .[torch]
- name: Update metadata
env:
HF_TOKEN: ${{ secrets.SAYAK_HF_TOKEN }}
run: |
python utils/update_metadata.py --commit_sha ${{ github.sha }}

View File

@@ -19,16 +19,6 @@ authors:
family-names: Rasul
- given-names: Mishig
family-names: Davaadorj
- given-names: Dhruv
family-names: Nair
- given-names: Sayak
family-names: Paul
- given-names: Steven
family-names: Liu
- given-names: William
family-names: Berman
- given-names: Yiyi
family-names: Xu
- given-names: Thomas
family-names: Wolf
repository-code: 'https://github.com/huggingface/diffusers'

View File

@@ -355,7 +355,7 @@ You will need basic `git` proficiency to be able to contribute to
manual. Type `git --help` in a shell and enjoy. If you prefer books, [Pro
Git](https://git-scm.com/book/en/v2) is a very good reference.
Follow these steps to start contributing ([supported Python versions](https://github.com/huggingface/diffusers/blob/42f25d601a910dceadaee6c44345896b4cfa9928/setup.py#L270)):
Follow these steps to start contributing ([supported Python versions](https://github.com/huggingface/diffusers/blob/main/setup.py#L265)):
1. Fork the [repository](https://github.com/huggingface/diffusers) by
clicking on the 'Fork' button on the repository's page. This creates a copy of the code

View File

@@ -42,7 +42,6 @@ repo-consistency:
quality:
ruff check $(check_dirs) setup.py
ruff format --check $(check_dirs) setup.py
doc-builder style src/diffusers docs/source --max_len 119 --check_only
python utils/check_doc_toc.py
# Format source code automatically and check is there are any problems left that need manual fixing
@@ -56,7 +55,6 @@ extra_style_checks:
style:
ruff check $(check_dirs) setup.py --fix
ruff format $(check_dirs) setup.py
doc-builder style src/diffusers docs/source --max_len 119
${MAKE} autogenerate_code
${MAKE} extra_style_checks

View File

@@ -77,7 +77,7 @@ Please refer to the [How to use Stable Diffusion in Apple Silicon](https://huggi
## Quickstart
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 25.000+ checkpoints):
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 19000+ checkpoints):
```python
from diffusers import DiffusionPipeline
@@ -219,7 +219,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
- https://github.com/deep-floyd/IF
- https://github.com/bentoml/BentoML
- https://github.com/bmaltais/kohya_ss
- +11.000 other amazing GitHub repositories 💪
- +8000 other amazing GitHub repositories 💪
Thank you for using us ❤️.
@@ -238,7 +238,7 @@ We also want to thank @heejkoo for the very helpful overview of papers, code and
```bibtex
@misc{von-platen-etal-2022-diffusers,
author = {Patrick von Platen and Suraj Patil and Anton Lozhkov and Pedro Cuenca and Nathan Lambert and Kashif Rasul and Mishig Davaadorj and Dhruv Nair and Sayak Paul and William Berman and Yiyi Xu and Steven Liu and Thomas Wolf},
author = {Patrick von Platen and Suraj Patil and Anton Lozhkov and Pedro Cuenca and Nathan Lambert and Kashif Rasul and Mishig Davaadorj and Thomas Wolf},
title = {Diffusers: State-of-the-art diffusion models},
year = {2022},
publisher = {GitHub},

View File

@@ -1,52 +0,0 @@
FROM ubuntu:20.04
LABEL maintainer="Hugging Face"
LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt-get -y update \
&& apt-get install -y software-properties-common \
&& add-apt-repository ppa:deadsnakes/ppa
RUN apt install -y bash \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.10 \
python3-pip \
libgl1 \
zip \
wget \
python3.10-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.10 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3.10 -m uv pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark \
--extra-index-url https://download.pytorch.org/whl/cpu && \
python3.10 -m uv pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy==1.26.4 \
scipy \
tensorboard \
transformers \
matplotlib \
setuptools==69.5.1
CMD ["/bin/bash"]

View File

@@ -4,25 +4,21 @@ LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt-get -y update \
&& apt-get install -y software-properties-common \
&& add-apt-repository ppa:deadsnakes/ppa
RUN apt install -y bash \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.10 \
python3-pip \
python3.10-venv && \
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.10 -m venv /opt/venv
RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
@@ -40,7 +36,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
huggingface-hub \
Jinja2 \
librosa \
numpy==1.26.4 \
numpy \
scipy \
tensorboard \
transformers

View File

@@ -4,25 +4,21 @@ LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt-get -y update \
&& apt-get install -y software-properties-common \
&& add-apt-repository ppa:deadsnakes/ppa
RUN apt install -y bash \
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.10 \
python3.8 \
python3-pip \
python3.10-venv && \
python3.8-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.10 -m venv /opt/venv
RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
@@ -41,8 +37,8 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy==1.26.4 \
librosa \
numpy \
scipy \
tensorboard \
transformers

View File

@@ -4,25 +4,21 @@ LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt-get -y update \
&& apt-get install -y software-properties-common \
&& add-apt-repository ppa:deadsnakes/ppa
RUN apt install -y bash \
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.10 \
python3.8 \
python3-pip \
python3.10-venv && \
python3.8-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.10 -m venv /opt/venv
RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
@@ -40,7 +36,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
huggingface-hub \
Jinja2 \
librosa \
numpy==1.26.4 \
numpy \
scipy \
tensorboard \
transformers

View File

@@ -4,43 +4,39 @@ LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt-get -y update \
&& apt-get install -y software-properties-common \
&& add-apt-repository ppa:deadsnakes/ppa
RUN apt install -y bash \
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.10 \
python3.8 \
python3-pip \
python3.10-venv && \
python3.8-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.10 -m venv /opt/venv
RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3.10 -m uv pip install --no-cache-dir \
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m uv pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
"onnxruntime-gpu>=1.13.1" \
--extra-index-url https://download.pytorch.org/whl/cu117 && \
python3.10 -m uv pip install --no-cache-dir \
python3 -m uv pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy==1.26.4 \
numpy \
scipy \
tensorboard \
transformers

View File

@@ -4,11 +4,8 @@ LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt-get -y update \
&& apt-get install -y software-properties-common \
&& add-apt-repository ppa:deadsnakes/ppa
RUN apt install -y bash \
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
@@ -16,30 +13,31 @@ RUN apt install -y bash \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.10 \
python3.9 \
python3.9-dev \
python3-pip \
python3.10-venv && \
python3.9-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.10 -m venv /opt/venv
RUN python3.9 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3.10 -m uv pip install --no-cache-dir \
RUN python3.9 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3.9 -m uv pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark && \
python3.10 -m pip install --no-cache-dir \
python3.9 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy==1.26.4 \
numpy \
scipy \
tensorboard \
transformers

View File

@@ -4,45 +4,42 @@ LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt-get -y update \
&& apt-get install -y software-properties-common \
&& add-apt-repository ppa:deadsnakes/ppa
RUN apt install -y bash \
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.10 \
python3.8 \
python3-pip \
libgl1 \
python3.10-venv && \
python3.8-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.10 -m venv /opt/venv
RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3.10 -m uv pip install --no-cache-dir \
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m uv pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark \
--extra-index-url https://download.pytorch.org/whl/cpu && \
python3.10 -m uv pip install --no-cache-dir \
python3 -m uv pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy==1.26.4 \
numpy \
scipy \
tensorboard \
transformers matplotlib
transformers
CMD ["/bin/bash"]

View File

@@ -4,11 +4,8 @@ LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt-get -y update \
&& apt-get install -y software-properties-common \
&& add-apt-repository ppa:deadsnakes/ppa
RUN apt install -y bash \
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
@@ -16,30 +13,30 @@ RUN apt install -y bash \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.10 \
python3.8 \
python3-pip \
python3.10-venv && \
python3.8-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.10 -m venv /opt/venv
RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3.10 -m uv pip install --no-cache-dir \
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m uv pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark && \
python3.10 -m pip install --no-cache-dir \
python3 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy==1.26.4 \
numpy \
scipy \
tensorboard \
transformers \

View File

@@ -4,11 +4,8 @@ LABEL repository="diffusers"
ENV DEBIAN_FRONTEND=noninteractive
RUN apt-get -y update \
&& apt-get install -y software-properties-common \
&& add-apt-repository ppa:deadsnakes/ppa
RUN apt install -y bash \
RUN apt update && \
apt install -y bash \
build-essential \
git \
git-lfs \
@@ -16,30 +13,30 @@ RUN apt install -y bash \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.10 \
python3.8 \
python3-pip \
python3.10-venv && \
python3.8-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.10 -m venv /opt/venv
RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3.10 -m pip install --no-cache-dir \
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark && \
python3.10 -m uv pip install --no-cache-dir \
python3 -m uv pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
numpy==1.26.4 \
numpy \
scipy \
tensorboard \
transformers \

View File

@@ -242,10 +242,10 @@ Here's an example of a tuple return, comprising several objects:
```
Returns:
`tuple(torch.Tensor)` comprising various elements depending on the configuration ([`BertConfig`]) and inputs:
- ** loss** (*optional*, returned when `masked_lm_labels` is provided) `torch.Tensor` of shape `(1,)` --
`tuple(torch.FloatTensor)` comprising various elements depending on the configuration ([`BertConfig`]) and inputs:
- ** loss** (*optional*, returned when `masked_lm_labels` is provided) `torch.FloatTensor` of shape `(1,)` --
Total loss is the sum of the masked language modeling loss and the next sequence prediction (classification) loss.
- **prediction_scores** (`torch.Tensor` of shape `(batch_size, sequence_length, config.vocab_size)`) --
- **prediction_scores** (`torch.FloatTensor` of shape `(batch_size, sequence_length, config.vocab_size)`) --
Prediction scores of the language modeling head (scores for each vocabulary token before SoftMax).
```

View File

@@ -18,145 +18,155 @@
- local: tutorials/basic_training
title: Train a diffusion model
- local: tutorials/using_peft_for_inference
title: Load LoRAs for inference
title: Inference with PEFT
- local: tutorials/fast_diffusion
title: Accelerate inference of text-to-image diffusion models
title: Tutorials
- sections:
- local: using-diffusers/loading
title: Load pipelines
- local: using-diffusers/custom_pipeline_overview
title: Load community pipelines and components
- local: using-diffusers/schedulers
title: Load schedulers and models
- local: using-diffusers/other-formats
title: Model files and layouts
- local: using-diffusers/loading_adapters
title: Load adapters
- local: using-diffusers/push_to_hub
title: Push files to the Hub
title: Load pipelines and adapters
- sections:
- local: using-diffusers/unconditional_image_generation
title: Unconditional image generation
- local: using-diffusers/conditional_image_generation
title: Text-to-image
- local: using-diffusers/img2img
title: Image-to-image
- local: using-diffusers/inpaint
title: Inpainting
- local: using-diffusers/text-img2vid
title: Text or image-to-video
- local: using-diffusers/depth2img
title: Depth-to-image
title: Generative tasks
- sections:
- local: using-diffusers/overview_techniques
title: Overview
- local: training/distributed_inference
title: Distributed inference with multiple GPUs
- local: using-diffusers/merge_loras
title: Merge LoRAs
- local: using-diffusers/scheduler_features
title: Scheduler features
- local: using-diffusers/callback
title: Pipeline callbacks
- local: using-diffusers/reusing_seeds
title: Reproducible pipelines
- local: using-diffusers/image_quality
title: Controlling image quality
- local: using-diffusers/weighted_prompts
title: Prompt techniques
title: Inference techniques
- sections:
- local: advanced_inference/outpaint
title: Outpainting
title: Advanced inference
- sections:
- local: using-diffusers/sdxl
title: Stable Diffusion XL
- local: using-diffusers/sdxl_turbo
title: SDXL Turbo
- local: using-diffusers/kandinsky
title: Kandinsky
- local: using-diffusers/ip_adapter
title: IP-Adapter
- local: using-diffusers/controlnet
title: ControlNet
- local: using-diffusers/t2i_adapter
title: T2I-Adapter
- local: using-diffusers/inference_with_lcm
title: Latent Consistency Model
- local: using-diffusers/textual_inversion_inference
title: Textual inversion
- local: using-diffusers/shap-e
title: Shap-E
- local: using-diffusers/diffedit
title: DiffEdit
- local: using-diffusers/inference_with_tcd_lora
title: Trajectory Consistency Distillation-LoRA
- local: using-diffusers/svd
title: Stable Video Diffusion
- local: using-diffusers/marigold_usage
title: Marigold Computer Vision
title: Specific pipeline examples
- sections:
- local: training/overview
title: Overview
- local: training/create_dataset
title: Create a dataset for training
- local: training/adapt_a_model
title: Adapt a model to a new task
- isExpanded: false
sections:
- local: training/unconditional_training
- sections:
- local: using-diffusers/loading_overview
title: Overview
- local: using-diffusers/loading
title: Load pipelines, models, and schedulers
- local: using-diffusers/schedulers
title: Load and compare different schedulers
- local: using-diffusers/custom_pipeline_overview
title: Load community pipelines and components
- local: using-diffusers/using_safetensors
title: Load safetensors
- local: using-diffusers/other-formats
title: Load different Stable Diffusion formats
- local: using-diffusers/loading_adapters
title: Load adapters
- local: using-diffusers/push_to_hub
title: Push files to the Hub
title: Loading & Hub
- sections:
- local: using-diffusers/pipeline_overview
title: Overview
- local: using-diffusers/unconditional_image_generation
title: Unconditional image generation
- local: training/text2image
- local: using-diffusers/conditional_image_generation
title: Text-to-image
- local: training/sdxl
- local: using-diffusers/img2img
title: Image-to-image
- local: using-diffusers/inpaint
title: Inpainting
- local: using-diffusers/text-img2vid
title: Text or image-to-video
- local: using-diffusers/depth2img
title: Depth-to-image
title: Tasks
- sections:
- local: using-diffusers/textual_inversion_inference
title: Textual inversion
- local: using-diffusers/ip_adapter
title: IP-Adapter
- local: training/distributed_inference
title: Distributed inference with multiple GPUs
- local: using-diffusers/reusing_seeds
title: Improve image quality with deterministic generation
- local: using-diffusers/control_brightness
title: Control image brightness
- local: using-diffusers/weighted_prompts
title: Prompt weighting
- local: using-diffusers/freeu
title: Improve generation quality with FreeU
title: Techniques
- sections:
- local: using-diffusers/pipeline_overview
title: Overview
- local: using-diffusers/sdxl
title: Stable Diffusion XL
- local: training/kandinsky
title: Kandinsky 2.2
- local: training/wuerstchen
title: Wuerstchen
- local: training/controlnet
- local: using-diffusers/sdxl_turbo
title: SDXL Turbo
- local: using-diffusers/kandinsky
title: Kandinsky
- local: using-diffusers/controlnet
title: ControlNet
- local: training/t2i_adapters
title: T2I-Adapters
- local: training/instructpix2pix
title: InstructPix2Pix
title: Models
- isExpanded: false
sections:
- local: training/text_inversion
title: Textual Inversion
- local: training/dreambooth
title: DreamBooth
- local: training/lora
title: LoRA
- local: training/custom_diffusion
title: Custom Diffusion
- local: training/lcm_distill
title: Latent Consistency Distillation
- local: training/ddpo
title: Reinforcement learning training with DDPO
title: Methods
title: Training
- local: using-diffusers/shap-e
title: Shap-E
- local: using-diffusers/diffedit
title: DiffEdit
- local: using-diffusers/distilled_sd
title: Distilled Stable Diffusion inference
- local: using-diffusers/callback
title: Pipeline callbacks
- local: using-diffusers/reproducibility
title: Create reproducible pipelines
- local: using-diffusers/custom_pipeline_examples
title: Community pipelines
- local: using-diffusers/contribute_pipeline
title: Contribute a community pipeline
- local: using-diffusers/inference_with_lcm_lora
title: Latent Consistency Model-LoRA
- local: using-diffusers/inference_with_lcm
title: Latent Consistency Model
- local: using-diffusers/svd
title: Stable Video Diffusion
title: Specific pipeline examples
- sections:
- local: training/overview
title: Overview
- local: training/create_dataset
title: Create a dataset for training
- local: training/adapt_a_model
title: Adapt a model to a new task
- sections:
- local: training/unconditional_training
title: Unconditional image generation
- local: training/text2image
title: Text-to-image
- local: training/sdxl
title: Stable Diffusion XL
- local: training/kandinsky
title: Kandinsky 2.2
- local: training/wuerstchen
title: Wuerstchen
- local: training/controlnet
title: ControlNet
- local: training/t2i_adapters
title: T2I-Adapters
- local: training/instructpix2pix
title: InstructPix2Pix
title: Models
- sections:
- local: training/text_inversion
title: Textual Inversion
- local: training/dreambooth
title: DreamBooth
- local: training/lora
title: LoRA
- local: training/custom_diffusion
title: Custom Diffusion
- local: training/lcm_distill
title: Latent Consistency Distillation
- local: training/ddpo
title: Reinforcement learning training with DDPO
title: Methods
title: Training
- sections:
- local: using-diffusers/other-modalities
title: Other Modalities
title: Taking Diffusers Beyond Images
title: Using Diffusers
- sections:
- local: optimization/fp16
title: Speed up inference
- local: optimization/memory
title: Reduce memory usage
- local: optimization/torch2.0
title: PyTorch 2.0
- local: optimization/xformers
title: xFormers
- local: optimization/tome
title: Token merging
- local: optimization/deepcache
title: DeepCache
- local: optimization/tgate
title: TGATE
- local: optimization/opt_overview
title: Overview
- sections:
- local: optimization/fp16
title: Speed up inference
- local: optimization/memory
title: Reduce memory usage
- local: optimization/torch2.0
title: PyTorch 2.0
- local: optimization/xformers
title: xFormers
- local: optimization/tome
title: Token merging
- local: optimization/deepcache
title: DeepCache
title: General optimizations
- sections:
- local: using-diffusers/stable_diffusion_jax_how_to
title: JAX/Flax
@@ -166,14 +176,14 @@
title: OpenVINO
- local: optimization/coreml
title: Core ML
title: Optimized model formats
title: Optimized model types
- sections:
- local: optimization/mps
title: Metal Performance Shaders (MPS)
- local: optimization/habana
title: Habana Gaudi
title: Optimized hardware
title: Accelerate inference and reduce memory
title: Optimization
- sections:
- local: conceptual/philosophy
title: Philosophy
@@ -187,8 +197,7 @@
title: Evaluating Diffusion Models
title: Conceptual Guides
- sections:
- isExpanded: false
sections:
- sections:
- local: api/configuration
title: Configuration
- local: api/logging
@@ -196,8 +205,7 @@
- local: api/outputs
title: Outputs
title: Main Classes
- isExpanded: false
sections:
- sections:
- local: api/loaders/ip_adapter
title: IP-Adapter
- local: api/loaders/lora
@@ -211,8 +219,7 @@
- local: api/loaders/peft
title: PEFT
title: Loaders
- isExpanded: false
sections:
- sections:
- local: api/models/overview
title: Overview
- local: api/models/unet
@@ -238,24 +245,15 @@
- local: api/models/consistency_decoder_vae
title: ConsistencyDecoderVAE
- local: api/models/transformer2d
title: Transformer2DModel
- local: api/models/pixart_transformer2d
title: PixArtTransformer2DModel
- local: api/models/dit_transformer2d
title: DiTTransformer2DModel
- local: api/models/hunyuan_transformer2d
title: HunyuanDiT2DModel
title: Transformer2D
- local: api/models/transformer_temporal
title: TransformerTemporalModel
- local: api/models/sd3_transformer2d
title: SD3Transformer2DModel
title: Transformer Temporal
- local: api/models/prior_transformer
title: PriorTransformer
title: Prior Transformer
- local: api/models/controlnet
title: ControlNetModel
title: ControlNet
title: Models
- isExpanded: false
sections:
- sections:
- local: api/pipelines/overview
title: Overview
- local: api/pipelines/amused
@@ -278,10 +276,6 @@
title: ControlNet
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/controlnetxs
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
title: ControlNet-XS with Stable Diffusion XL
- local: api/pipelines/dance_diffusion
title: Dance Diffusion
- local: api/pipelines/ddim
@@ -294,8 +288,6 @@
title: DiffEdit
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/hunyuandit
title: Hunyuan-DiT
- local: api/pipelines/i2vgenxl
title: I2VGen-XL
- local: api/pipelines/pix2pix
@@ -310,10 +302,6 @@
title: Latent Consistency Models
- local: api/pipelines/latent_diffusion
title: Latent Diffusion
- local: api/pipelines/ledits_pp
title: LEDITS++
- local: api/pipelines/marigold
title: Marigold
- local: api/pipelines/panorama
title: MultiDiffusion
- local: api/pipelines/musicldm
@@ -324,8 +312,6 @@
title: Personalized Image Animator (PIA)
- local: api/pipelines/pixart
title: PixArt-α
- local: api/pipelines/pixart_sigma
title: PixArt-Σ
- local: api/pipelines/self_attention_guidance
title: Self-Attention Guidance
- local: api/pipelines/semantic_stable_diffusion
@@ -353,8 +339,6 @@
title: Safe Stable Diffusion
- local: api/pipelines/stable_diffusion/stable_diffusion_2
title: Stable Diffusion 2
- local: api/pipelines/stable_diffusion/stable_diffusion_3
title: Stable Diffusion 3
- local: api/pipelines/stable_diffusion/stable_diffusion_xl
title: Stable Diffusion XL
- local: api/pipelines/stable_diffusion/sdxl_turbo
@@ -368,7 +352,7 @@
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D Upscaler
- local: api/pipelines/stable_diffusion/adapter
title: T2I-Adapter
title: Stable Diffusion T2I-Adapter
- local: api/pipelines/stable_diffusion/gligen
title: GLIGEN (Grounded Language-to-Image Generation)
title: Stable Diffusion
@@ -387,8 +371,7 @@
- local: api/pipelines/wuerstchen
title: Wuerstchen
title: Pipelines
- isExpanded: false
sections:
- sections:
- local: api/schedulers/overview
title: Overview
- local: api/schedulers/cm_stochastic_iterative
@@ -411,16 +394,10 @@
title: DPMSolverSDEScheduler
- local: api/schedulers/singlestep_dpm_solver
title: DPMSolverSinglestepScheduler
- local: api/schedulers/edm_multistep_dpm_solver
title: EDMDPMSolverMultistepScheduler
- local: api/schedulers/edm_euler
title: EDMEulerScheduler
- local: api/schedulers/euler_ancestral
title: EulerAncestralDiscreteScheduler
- local: api/schedulers/euler
title: EulerDiscreteScheduler
- local: api/schedulers/flow_match_euler_discrete
title: FlowMatchEulerDiscreteScheduler
- local: api/schedulers/heun
title: HeunDiscreteScheduler
- local: api/schedulers/ipndm
@@ -450,8 +427,7 @@
- local: api/schedulers/vq_diffusion
title: VQDiffusionScheduler
title: Schedulers
- isExpanded: false
sections:
- sections:
- local: api/internal_classes_overview
title: Overview
- local: api/attnprocessor
@@ -464,7 +440,5 @@
title: Utilities
- local: api/image_processor
title: VAE Image Processor
- local: api/video_processor
title: Video Processor
title: Internal classes
title: API

View File

@@ -1,231 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Outpainting
Outpainting extends an image beyond its original boundaries, allowing you to add, replace, or modify visual elements in an image while preserving the original image. Like [inpainting](../using-diffusers/inpaint), you want to fill the white area (in this case, the area outside of the original image) with new visual elements while keeping the original image (represented by a mask of black pixels). There are a couple of ways to outpaint, such as with a [ControlNet](https://hf.co/blog/OzzyGT/outpainting-controlnet) or with [Differential Diffusion](https://hf.co/blog/OzzyGT/outpainting-differential-diffusion).
This guide will show you how to outpaint with an inpainting model, ControlNet, and a ZoeDepth estimator.
Before you begin, make sure you have the [controlnet_aux](https://github.com/huggingface/controlnet_aux) library installed so you can use the ZoeDepth estimator.
```py
!pip install -q controlnet_aux
```
## Image preparation
Start by picking an image to outpaint with and remove the background with a Space like [BRIA-RMBG-1.4](https://hf.co/spaces/briaai/BRIA-RMBG-1.4).
<iframe
src="https://briaai-bria-rmbg-1-4.hf.space"
frameborder="0"
width="850"
height="450"
></iframe>
For example, remove the background from this image of a pair of shoes.
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/stevhliu/testing-images/resolve/main/original-jordan.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">original image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/stevhliu/testing-images/resolve/main/no-background-jordan.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">background removed</figcaption>
</div>
</div>
[Stable Diffusion XL (SDXL)](../using-diffusers/sdxl) models work best with 1024x1024 images, but you can resize the image to any size as long as your hardware has enough memory to support it. The transparent background in the image should also be replaced with a white background. Create a function (like the one below) that scales and pastes the image onto a white background.
```py
import random
import requests
import torch
from controlnet_aux import ZoeDetector
from PIL import Image, ImageOps
from diffusers import (
AutoencoderKL,
ControlNetModel,
StableDiffusionXLControlNetPipeline,
StableDiffusionXLInpaintPipeline,
)
def scale_and_paste(original_image):
aspect_ratio = original_image.width / original_image.height
if original_image.width > original_image.height:
new_width = 1024
new_height = round(new_width / aspect_ratio)
else:
new_height = 1024
new_width = round(new_height * aspect_ratio)
resized_original = original_image.resize((new_width, new_height), Image.LANCZOS)
white_background = Image.new("RGBA", (1024, 1024), "white")
x = (1024 - new_width) // 2
y = (1024 - new_height) // 2
white_background.paste(resized_original, (x, y), resized_original)
return resized_original, white_background
original_image = Image.open(
requests.get(
"https://huggingface.co/datasets/stevhliu/testing-images/resolve/main/no-background-jordan.png",
stream=True,
).raw
).convert("RGBA")
resized_img, white_bg_image = scale_and_paste(original_image)
```
To avoid adding unwanted extra details, use the ZoeDepth estimator to provide additional guidance during generation and to ensure the shoes remain consistent with the original image.
```py
zoe = ZoeDetector.from_pretrained("lllyasviel/Annotators")
image_zoe = zoe(white_bg_image, detect_resolution=512, image_resolution=1024)
image_zoe
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/stevhliu/testing-images/resolve/main/zoedepth-jordan.png"/>
</div>
## Outpaint
Once your image is ready, you can generate content in the white area around the shoes with [controlnet-inpaint-dreamer-sdxl](https://hf.co/destitech/controlnet-inpaint-dreamer-sdxl), a SDXL ControlNet trained for inpainting.
Load the inpainting ControlNet, ZoeDepth model, VAE and pass them to the [`StableDiffusionXLControlNetPipeline`]. Then you can create an optional `generate_image` function (for convenience) to outpaint an initial image.
```py
controlnets = [
ControlNetModel.from_pretrained(
"destitech/controlnet-inpaint-dreamer-sdxl", torch_dtype=torch.float16, variant="fp16"
),
ControlNetModel.from_pretrained(
"diffusers/controlnet-zoe-depth-sdxl-1.0", torch_dtype=torch.float16
),
]
vae = AutoencoderKL.from_pretrained("madebyollin/sdxl-vae-fp16-fix", torch_dtype=torch.float16).to("cuda")
pipeline = StableDiffusionXLControlNetPipeline.from_pretrained(
"SG161222/RealVisXL_V4.0", torch_dtype=torch.float16, variant="fp16", controlnet=controlnets, vae=vae
).to("cuda")
def generate_image(prompt, negative_prompt, inpaint_image, zoe_image, seed: int = None):
if seed is None:
seed = random.randint(0, 2**32 - 1)
generator = torch.Generator(device="cpu").manual_seed(seed)
image = pipeline(
prompt,
negative_prompt=negative_prompt,
image=[inpaint_image, zoe_image],
guidance_scale=6.5,
num_inference_steps=25,
generator=generator,
controlnet_conditioning_scale=[0.5, 0.8],
control_guidance_end=[0.9, 0.6],
).images[0]
return image
prompt = "nike air jordans on a basketball court"
negative_prompt = ""
temp_image = generate_image(prompt, negative_prompt, white_bg_image, image_zoe, 908097)
```
Paste the original image over the initial outpainted image. You'll improve the outpainted background in a later step.
```py
x = (1024 - resized_img.width) // 2
y = (1024 - resized_img.height) // 2
temp_image.paste(resized_img, (x, y), resized_img)
temp_image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/stevhliu/testing-images/resolve/main/initial-outpaint.png"/>
</div>
> [!TIP]
> Now is a good time to free up some memory if you're running low!
>
> ```py
> pipeline=None
> torch.cuda.empty_cache()
> ```
Now that you have an initial outpainted image, load the [`StableDiffusionXLInpaintPipeline`] with the [RealVisXL](https://hf.co/SG161222/RealVisXL_V4.0) model to generate the final outpainted image with better quality.
```py
pipeline = StableDiffusionXLInpaintPipeline.from_pretrained(
"OzzyGT/RealVisXL_V4.0_inpainting",
torch_dtype=torch.float16,
variant="fp16",
vae=vae,
).to("cuda")
```
Prepare a mask for the final outpainted image. To create a more natural transition between the original image and the outpainted background, blur the mask to help it blend better.
```py
mask = Image.new("L", temp_image.size)
mask.paste(resized_img.split()[3], (x, y))
mask = ImageOps.invert(mask)
final_mask = mask.point(lambda p: p > 128 and 255)
mask_blurred = pipeline.mask_processor.blur(final_mask, blur_factor=20)
mask_blurred
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/stevhliu/testing-images/resolve/main/blurred-mask.png"/>
</div>
Create a better prompt and pass it to the `generate_outpaint` function to generate the final outpainted image. Again, paste the original image over the final outpainted background.
```py
def generate_outpaint(prompt, negative_prompt, image, mask, seed: int = None):
if seed is None:
seed = random.randint(0, 2**32 - 1)
generator = torch.Generator(device="cpu").manual_seed(seed)
image = pipeline(
prompt,
negative_prompt=negative_prompt,
image=image,
mask_image=mask,
guidance_scale=10.0,
strength=0.8,
num_inference_steps=30,
generator=generator,
).images[0]
return image
prompt = "high quality photo of nike air jordans on a basketball court, highly detailed"
negative_prompt = ""
final_image = generate_outpaint(prompt, negative_prompt, temp_image, mask_blurred, 7688778)
x = (1024 - resized_img.width) // 2
y = (1024 - resized_img.height) // 2
final_image.paste(resized_img, (x, y), resized_img)
final_image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/stevhliu/testing-images/resolve/main/final-outpaint.png"/>
</div>

View File

@@ -55,6 +55,3 @@ An attention processor is a class for applying different types of attention mech
## XFormersAttnProcessor
[[autodoc]] models.attention_processor.XFormersAttnProcessor
## AttnProcessorNPU
[[autodoc]] models.attention_processor.AttnProcessorNPU

View File

@@ -25,11 +25,3 @@ All pipelines with [`VaeImageProcessor`] accept PIL Image, PyTorch tensor, or Nu
The [`VaeImageProcessorLDM3D`] accepts RGB and depth inputs and returns RGB and depth outputs.
[[autodoc]] image_processor.VaeImageProcessorLDM3D
## PixArtImageProcessor
[[autodoc]] image_processor.PixArtImageProcessor
## IPAdapterMaskProcessor
[[autodoc]] image_processor.IPAdapterMaskProcessor

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@@ -23,7 +23,3 @@ Learn how to load an IP-Adapter checkpoint and image in the IP-Adapter [loading]
## IPAdapterMixin
[[autodoc]] loaders.ip_adapter.IPAdapterMixin
## IPAdapterMaskProcessor
[[autodoc]] image_processor.IPAdapterMaskProcessor

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@@ -12,50 +12,26 @@ specific language governing permissions and limitations under the License.
# Single files
The [`~loaders.FromSingleFileMixin.from_single_file`] method allows you to load:
Diffusers supports loading pretrained pipeline (or model) weights stored in a single file, such as a `ckpt` or `safetensors` file. These single file types are typically produced from community trained models. There are three classes for loading single file weights:
* a model stored in a single file, which is useful if you're working with models from the diffusion ecosystem, like Automatic1111, and commonly rely on a single-file layout to store and share models
* a model stored in their originally distributed layout, which is useful if you're working with models finetuned with other services, and want to load it directly into Diffusers model objects and pipelines
- [`FromSingleFileMixin`] supports loading pretrained pipeline weights stored in a single file, which can either be a `ckpt` or `safetensors` file.
- [`FromOriginalVAEMixin`] supports loading a pretrained [`AutoencoderKL`] from pretrained ControlNet weights stored in a single file, which can either be a `ckpt` or `safetensors` file.
- [`FromOriginalControlnetMixin`] supports loading pretrained ControlNet weights stored in a single file, which can either be a `ckpt` or `safetensors` file.
> [!TIP]
> Read the [Model files and layouts](../../using-diffusers/other-formats) guide to learn more about the Diffusers-multifolder layout versus the single-file layout, and how to load models stored in these different layouts.
<Tip>
## Supported pipelines
To learn more about how to load single file weights, see the [Load different Stable Diffusion formats](../../using-diffusers/other-formats) loading guide.
- [`StableDiffusionPipeline`]
- [`StableDiffusionImg2ImgPipeline`]
- [`StableDiffusionInpaintPipeline`]
- [`StableDiffusionControlNetPipeline`]
- [`StableDiffusionControlNetImg2ImgPipeline`]
- [`StableDiffusionControlNetInpaintPipeline`]
- [`StableDiffusionUpscalePipeline`]
- [`StableDiffusionXLPipeline`]
- [`StableDiffusionXLImg2ImgPipeline`]
- [`StableDiffusionXLInpaintPipeline`]
- [`StableDiffusionXLInstructPix2PixPipeline`]
- [`StableDiffusionXLControlNetPipeline`]
- [`StableDiffusionXLKDiffusionPipeline`]
- [`StableDiffusion3Pipeline`]
- [`LatentConsistencyModelPipeline`]
- [`LatentConsistencyModelImg2ImgPipeline`]
- [`StableDiffusionControlNetXSPipeline`]
- [`StableDiffusionXLControlNetXSPipeline`]
- [`LEditsPPPipelineStableDiffusion`]
- [`LEditsPPPipelineStableDiffusionXL`]
- [`PIAPipeline`]
## Supported models
- [`UNet2DConditionModel`]
- [`StableCascadeUNet`]
- [`AutoencoderKL`]
- [`ControlNetModel`]
- [`SD3Transformer2DModel`]
</Tip>
## FromSingleFileMixin
[[autodoc]] loaders.single_file.FromSingleFileMixin
## FromOriginalModelMixin
## FromOriginalVAEMixin
[[autodoc]] loaders.single_file_model.FromOriginalModelMixin
[[autodoc]] loaders.autoencoder.FromOriginalVAEMixin
## FromOriginalControlnetMixin
[[autodoc]] loaders.controlnet.FromOriginalControlNetMixin

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@@ -10,7 +10,7 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# ControlNetModel
# ControlNet
The ControlNet model was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, Maneesh Agrawala. It provides a greater degree of control over text-to-image generation by conditioning the model on additional inputs such as edge maps, depth maps, segmentation maps, and keypoints for pose detection.

View File

@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# DiTTransformer2DModel
A Transformer model for image-like data from [DiT](https://huggingface.co/papers/2212.09748).
## DiTTransformer2DModel
[[autodoc]] DiTTransformer2DModel

View File

@@ -1,20 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# HunyuanDiT2DModel
A Diffusion Transformer model for 2D data from [Hunyuan-DiT](https://github.com/Tencent/HunyuanDiT).
## HunyuanDiT2DModel
[[autodoc]] HunyuanDiT2DModel

View File

@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# PixArtTransformer2DModel
A Transformer model for image-like data from [PixArt-Alpha](https://huggingface.co/papers/2310.00426) and [PixArt-Sigma](https://huggingface.co/papers/2403.04692).
## PixArtTransformer2DModel
[[autodoc]] PixArtTransformer2DModel

View File

@@ -10,7 +10,7 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# PriorTransformer
# Prior Transformer
The Prior Transformer was originally introduced in [Hierarchical Text-Conditional Image Generation with CLIP Latents](https://huggingface.co/papers/2204.06125) by Ramesh et al. It is used to predict CLIP image embeddings from CLIP text embeddings; image embeddings are predicted through a denoising diffusion process.

View File

@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# SD3 Transformer Model
The Transformer model introduced in [Stable Diffusion 3](https://hf.co/papers/2403.03206). Its novelty lies in the MMDiT transformer block.
## SD3Transformer2DModel
[[autodoc]] SD3Transformer2DModel

View File

@@ -10,7 +10,7 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Transformer2DModel
# Transformer2D
A Transformer model for image-like data from [CompVis](https://huggingface.co/CompVis) that is based on the [Vision Transformer](https://huggingface.co/papers/2010.11929) introduced by Dosovitskiy et al. The [`Transformer2DModel`] accepts discrete (classes of vector embeddings) or continuous (actual embeddings) inputs.

View File

@@ -10,7 +10,7 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# TransformerTemporalModel
# Transformer Temporal
A Transformer model for video-like data.

View File

@@ -24,4 +24,4 @@ The abstract from the paper is:
## VQEncoderOutput
[[autodoc]] models.autoencoders.vq_model.VQEncoderOutput
[[autodoc]] models.vq_model.VQEncoderOutput

View File

@@ -16,7 +16,7 @@ aMUSEd was introduced in [aMUSEd: An Open MUSE Reproduction](https://huggingface
Amused is a lightweight text to image model based off of the [MUSE](https://arxiv.org/abs/2301.00704) architecture. Amused is particularly useful in applications that require a lightweight and fast model such as generating many images quickly at once.
Amused is a vqvae token based transformer that can generate an image in fewer forward passes than many diffusion models. In contrast with muse, it uses the smaller text encoder CLIP-L/14 instead of t5-xxl. Due to its small parameter count and few forward pass generation process, amused can generate many images quickly. This benefit is seen particularly at larger batch sizes.
Amused is a vqvae token based transformer that can generate an image in fewer forward passes than many diffusion models. In contrast with muse, it uses the smaller text encoder CLIP-L/14 instead of t5-xxl. Due to its small parameter count and few forward pass generation process, amused can generate many images quickly. This benefit is seen particularly at larger batch sizes.
The abstract from the paper is:

View File

@@ -101,53 +101,6 @@ AnimateDiff tends to work better with finetuned Stable Diffusion models. If you
</Tip>
### AnimateDiffSDXLPipeline
AnimateDiff can also be used with SDXL models. This is currently an experimental feature as only a beta release of the motion adapter checkpoint is available.
```python
import torch
from diffusers.models import MotionAdapter
from diffusers import AnimateDiffSDXLPipeline, DDIMScheduler
from diffusers.utils import export_to_gif
adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-sdxl-beta", torch_dtype=torch.float16)
model_id = "stabilityai/stable-diffusion-xl-base-1.0"
scheduler = DDIMScheduler.from_pretrained(
model_id,
subfolder="scheduler",
clip_sample=False,
timestep_spacing="linspace",
beta_schedule="linear",
steps_offset=1,
)
pipe = AnimateDiffSDXLPipeline.from_pretrained(
model_id,
motion_adapter=adapter,
scheduler=scheduler,
torch_dtype=torch.float16,
variant="fp16",
).to("cuda")
# enable memory savings
pipe.enable_vae_slicing()
pipe.enable_vae_tiling()
output = pipe(
prompt="a panda surfing in the ocean, realistic, high quality",
negative_prompt="low quality, worst quality",
num_inference_steps=20,
guidance_scale=8,
width=1024,
height=1024,
num_frames=16,
)
frames = output.frames[0]
export_to_gif(frames, "animation.gif")
```
### AnimateDiffVideoToVideoPipeline
AnimateDiff can also be used to generate visually similar videos or enable style/character/background or other edits starting from an initial video, allowing you to seamlessly explore creative possibilities.
@@ -165,7 +118,7 @@ from PIL import Image
adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-v1-5-2", torch_dtype=torch.float16)
# load SD 1.5 based finetuned model
model_id = "SG161222/Realistic_Vision_V5.1_noVAE"
pipe = AnimateDiffVideoToVideoPipeline.from_pretrained(model_id, motion_adapter=adapter, torch_dtype=torch.float16)
pipe = AnimateDiffVideoToVideoPipeline.from_pretrained(model_id, motion_adapter=adapter, torch_dtype=torch.float16).to("cuda")
scheduler = DDIMScheduler.from_pretrained(
model_id,
subfolder="scheduler",
@@ -455,29 +408,6 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
</Tip>
<table>
<tr>
<th align=center>Without FreeInit enabled</th>
<th align=center>With FreeInit enabled</th>
</tr>
<tr>
<td align=center>
panda playing a guitar
<br />
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-no-freeinit.gif"
alt="panda playing a guitar"
style="width: 300px;" />
</td>
<td align=center>
panda playing a guitar
<br/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-freeinit.gif"
alt="panda playing a guitar"
style="width: 300px;" />
</td>
</tr>
</table>
## Using AnimateLCM
[AnimateLCM](https://animatelcm.github.io/) is a motion module checkpoint and an [LCM LoRA](https://huggingface.co/docs/diffusers/using-diffusers/inference_with_lcm_lora) that have been created using a consistency learning strategy that decouples the distillation of the image generation priors and the motion generation priors.
@@ -569,12 +499,6 @@ export_to_gif(frames, "animatelcm-motion-lora.gif")
- all
- __call__
## AnimateDiffSDXLPipeline
[[autodoc]] AnimateDiffSDXLPipeline
- all
- __call__
## AnimateDiffVideoToVideoPipeline
[[autodoc]] AnimateDiffVideoToVideoPipeline

View File

@@ -20,8 +20,7 @@ The abstract of the paper is the following:
*Although audio generation shares commonalities across different types of audio, such as speech, music, and sound effects, designing models for each type requires careful consideration of specific objectives and biases that can significantly differ from those of other types. To bring us closer to a unified perspective of audio generation, this paper proposes a framework that utilizes the same learning method for speech, music, and sound effect generation. Our framework introduces a general representation of audio, called "language of audio" (LOA). Any audio can be translated into LOA based on AudioMAE, a self-supervised pre-trained representation learning model. In the generation process, we translate any modalities into LOA by using a GPT-2 model, and we perform self-supervised audio generation learning with a latent diffusion model conditioned on LOA. The proposed framework naturally brings advantages such as in-context learning abilities and reusable self-supervised pretrained AudioMAE and latent diffusion models. Experiments on the major benchmarks of text-to-audio, text-to-music, and text-to-speech demonstrate state-of-the-art or competitive performance against previous approaches. Our code, pretrained model, and demo are available at [this https URL](https://audioldm.github.io/audioldm2).*
This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi) and [Nguyễn Công Tú Anh](https://github.com/tuanh123789). The original codebase can be
found at [haoheliu/audioldm2](https://github.com/haoheliu/audioldm2).
This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi). The original codebase can be found at [haoheliu/audioldm2](https://github.com/haoheliu/audioldm2).
## Tips
@@ -37,8 +36,6 @@ See table below for details on the three checkpoints:
| [audioldm2](https://huggingface.co/cvssp/audioldm2) | Text-to-audio | 350M | 1.1B | 1150k |
| [audioldm2-large](https://huggingface.co/cvssp/audioldm2-large) | Text-to-audio | 750M | 1.5B | 1150k |
| [audioldm2-music](https://huggingface.co/cvssp/audioldm2-music) | Text-to-music | 350M | 1.1B | 665k |
| [audioldm2-gigaspeech](https://huggingface.co/anhnct/audioldm2_gigaspeech) | Text-to-speech | 350M | 1.1B |10k |
| [audioldm2-ljspeech](https://huggingface.co/anhnct/audioldm2_ljspeech) | Text-to-speech | 350M | 1.1B | |
### Constructing a prompt
@@ -56,7 +53,7 @@ See table below for details on the three checkpoints:
* The quality of the generated waveforms can vary significantly based on the seed. Try generating with different seeds until you find a satisfactory generation.
* Multiple waveforms can be generated in one go: set `num_waveforms_per_prompt` to a value greater than 1. Automatic scoring will be performed between the generated waveforms and prompt text, and the audios ranked from best to worst accordingly.
The following example demonstrates how to construct good music and speech generation using the aforementioned tips: [example](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2Pipeline.__call__.example).
The following example demonstrates how to construct good music generation using the aforementioned tips: [example](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2Pipeline.__call__.example).
<Tip>

View File

@@ -12,10 +12,42 @@ specific language governing permissions and limitations under the License.
# AutoPipeline
The `AutoPipeline` is designed to make it easy to load a checkpoint for a task without needing to know the specific pipeline class. Based on the task, the `AutoPipeline` automatically retrieves the correct pipeline class from the checkpoint `model_index.json` file.
`AutoPipeline` is designed to:
1. make it easy for you to load a checkpoint for a task without knowing the specific pipeline class to use
2. use multiple pipelines in your workflow
Based on the task, the `AutoPipeline` class automatically retrieves the relevant pipeline given the name or path to the pretrained weights with the `from_pretrained()` method.
To seamlessly switch between tasks with the same checkpoint without reallocating additional memory, use the `from_pipe()` method to transfer the components from the original pipeline to the new one.
```py
from diffusers import AutoPipelineForText2Image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
image = pipeline(prompt, num_inference_steps=25).images[0]
```
<Tip>
Check out the [AutoPipeline](../../tutorials/autopipeline) tutorial to learn how to use this API!
</Tip>
`AutoPipeline` supports text-to-image, image-to-image, and inpainting for the following diffusion models:
- [Stable Diffusion](./stable_diffusion/overview)
- [ControlNet](./controlnet)
- [Stable Diffusion XL (SDXL)](./stable_diffusion/stable_diffusion_xl)
- [DeepFloyd IF](./deepfloyd_if)
- [Kandinsky 2.1](./kandinsky)
- [Kandinsky 2.2](./kandinsky_v22)
> [!TIP]
> Check out the [AutoPipeline](../../tutorials/autopipeline) tutorial to learn how to use this API!
## AutoPipelineForText2Image

View File

@@ -1,95 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Hunyuan-DiT
![chinese elements understanding](https://github.com/gnobitab/diffusers-hunyuan/assets/1157982/39b99036-c3cb-4f16-bb1a-40ec25eda573)
[Hunyuan-DiT : A Powerful Multi-Resolution Diffusion Transformer with Fine-Grained Chinese Understanding](https://arxiv.org/abs/2405.08748) from Tencent Hunyuan.
The abstract from the paper is:
*We present Hunyuan-DiT, a text-to-image diffusion transformer with fine-grained understanding of both English and Chinese. To construct Hunyuan-DiT, we carefully design the transformer structure, text encoder, and positional encoding. We also build from scratch a whole data pipeline to update and evaluate data for iterative model optimization. For fine-grained language understanding, we train a Multimodal Large Language Model to refine the captions of the images. Finally, Hunyuan-DiT can perform multi-turn multimodal dialogue with users, generating and refining images according to the context. Through our holistic human evaluation protocol with more than 50 professional human evaluators, Hunyuan-DiT sets a new state-of-the-art in Chinese-to-image generation compared with other open-source models.*
You can find the original codebase at [Tencent/HunyuanDiT](https://github.com/Tencent/HunyuanDiT) and all the available checkpoints at [Tencent-Hunyuan](https://huggingface.co/Tencent-Hunyuan/HunyuanDiT).
**Highlights**: HunyuanDiT supports Chinese/English-to-image, multi-resolution generation.
HunyuanDiT has the following components:
* It uses a diffusion transformer as the backbone
* It combines two text encoders, a bilingual CLIP and a multilingual T5 encoder
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Optimization
You can optimize the pipeline's runtime and memory consumption with torch.compile and feed-forward chunking. To learn about other optimization methods, check out the [Speed up inference](../../optimization/fp16) and [Reduce memory usage](../../optimization/memory) guides.
### Inference
Use [`torch.compile`](https://huggingface.co/docs/diffusers/main/en/tutorials/fast_diffusion#torchcompile) to reduce the inference latency.
First, load the pipeline:
```python
from diffusers import HunyuanDiTPipeline
import torch
pipeline = HunyuanDiTPipeline.from_pretrained(
"Tencent-Hunyuan/HunyuanDiT-Diffusers", torch_dtype=torch.float16
).to("cuda")
```
Then change the memory layout of the pipelines `transformer` and `vae` components to `torch.channels-last`:
```python
pipeline.transformer.to(memory_format=torch.channels_last)
pipeline.vae.to(memory_format=torch.channels_last)
```
Finally, compile the components and run inference:
```python
pipeline.transformer = torch.compile(pipeline.transformer, mode="max-autotune", fullgraph=True)
pipeline.vae.decode = torch.compile(pipeline.vae.decode, mode="max-autotune", fullgraph=True)
image = pipeline(prompt="一个宇航员在骑马").images[0]
```
The [benchmark](https://gist.github.com/sayakpaul/29d3a14905cfcbf611fe71ebd22e9b23) results on a 80GB A100 machine are:
```bash
With torch.compile(): Average inference time: 12.470 seconds.
Without torch.compile(): Average inference time: 20.570 seconds.
```
### Memory optimization
By loading the T5 text encoder in 8 bits, you can run the pipeline in just under 6 GBs of GPU VRAM. Refer to [this script](https://gist.github.com/sayakpaul/3154605f6af05b98a41081aaba5ca43e) for details.
Furthermore, you can use the [`~HunyuanDiT2DModel.enable_forward_chunking`] method to reduce memory usage. Feed-forward chunking runs the feed-forward layers in a transformer block in a loop instead of all at once. This gives you a trade-off between memory consumption and inference runtime.
```diff
+ pipeline.transformer.enable_forward_chunking(chunk_size=1, dim=1)
```
## HunyuanDiTPipeline
[[autodoc]] HunyuanDiTPipeline
- all
- __call__

View File

@@ -47,7 +47,6 @@ Sample output with I2VGenXL:
* Unlike SVD, it additionally accepts text prompts as inputs.
* It can generate higher resolution videos.
* When using the [`DDIMScheduler`] (which is default for this pipeline), less than 50 steps for inference leads to bad results.
* This implementation is 1-stage variant of I2VGenXL. The main figure in the [I2VGen-XL](https://arxiv.org/abs/2311.04145) paper shows a 2-stage variant, however, 1-stage variant works well. See [this discussion](https://github.com/huggingface/diffusers/discussions/7952) for more details.
## I2VGenXLPipeline
[[autodoc]] I2VGenXLPipeline

View File

@@ -11,12 +11,12 @@ specific language governing permissions and limitations under the License.
Kandinsky 3 is created by [Vladimir Arkhipkin](https://github.com/oriBetelgeuse),[Anastasia Maltseva](https://github.com/NastyaMittseva),[Igor Pavlov](https://github.com/boomb0om),[Andrei Filatov](https://github.com/anvilarth),[Arseniy Shakhmatov](https://github.com/cene555),[Andrey Kuznetsov](https://github.com/kuznetsoffandrey),[Denis Dimitrov](https://github.com/denndimitrov), [Zein Shaheen](https://github.com/zeinsh)
The description from it's Github page:
The description from it's Github page:
*Kandinsky 3.0 is an open-source text-to-image diffusion model built upon the Kandinsky2-x model family. In comparison to its predecessors, enhancements have been made to the text understanding and visual quality of the model, achieved by increasing the size of the text encoder and Diffusion U-Net models, respectively.*
Its architecture includes 3 main components:
1. [FLAN-UL2](https://huggingface.co/google/flan-ul2), which is an encoder decoder model based on the T5 architecture.
1. [FLAN-UL2](https://huggingface.co/google/flan-ul2), which is an encoder decoder model based on the T5 architecture.
2. New U-Net architecture featuring BigGAN-deep blocks doubles depth while maintaining the same number of parameters.
3. Sber-MoVQGAN is a decoder proven to have superior results in image restoration.

View File

@@ -1,54 +0,0 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# LEDITS++
LEDITS++ was proposed in [LEDITS++: Limitless Image Editing using Text-to-Image Models](https://huggingface.co/papers/2311.16711) by Manuel Brack, Felix Friedrich, Katharina Kornmeier, Linoy Tsaban, Patrick Schramowski, Kristian Kersting, Apolinário Passos.
The abstract from the paper is:
*Text-to-image diffusion models have recently received increasing interest for their astonishing ability to produce high-fidelity images from solely text inputs. Subsequent research efforts aim to exploit and apply their capabilities to real image editing. However, existing image-to-image methods are often inefficient, imprecise, and of limited versatility. They either require time-consuming fine-tuning, deviate unnecessarily strongly from the input image, and/or lack support for multiple, simultaneous edits. To address these issues, we introduce LEDITS++, an efficient yet versatile and precise textual image manipulation technique. LEDITS++'s novel inversion approach requires no tuning nor optimization and produces high-fidelity results with a few diffusion steps. Second, our methodology supports multiple simultaneous edits and is architecture-agnostic. Third, we use a novel implicit masking technique that limits changes to relevant image regions. We propose the novel TEdBench++ benchmark as part of our exhaustive evaluation. Our results demonstrate the capabilities of LEDITS++ and its improvements over previous methods. The project page is available at https://leditsplusplus-project.static.hf.space .*
<Tip>
You can find additional information about LEDITS++ on the [project page](https://leditsplusplus-project.static.hf.space/index.html) and try it out in a [demo](https://huggingface.co/spaces/editing-images/leditsplusplus).
</Tip>
<Tip warning={true}>
Due to some backward compatability issues with the current diffusers implementation of [`~schedulers.DPMSolverMultistepScheduler`] this implementation of LEdits++ can no longer guarantee perfect inversion.
This issue is unlikely to have any noticeable effects on applied use-cases. However, we provide an alternative implementation that guarantees perfect inversion in a dedicated [GitHub repo](https://github.com/ml-research/ledits_pp).
</Tip>
We provide two distinct pipelines based on different pre-trained models.
## LEditsPPPipelineStableDiffusion
[[autodoc]] pipelines.ledits_pp.LEditsPPPipelineStableDiffusion
- all
- __call__
- invert
## LEditsPPPipelineStableDiffusionXL
[[autodoc]] pipelines.ledits_pp.LEditsPPPipelineStableDiffusionXL
- all
- __call__
- invert
## LEditsPPDiffusionPipelineOutput
[[autodoc]] pipelines.ledits_pp.pipeline_output.LEditsPPDiffusionPipelineOutput
- all
## LEditsPPInversionPipelineOutput
[[autodoc]] pipelines.ledits_pp.pipeline_output.LEditsPPInversionPipelineOutput
- all

View File

@@ -1,76 +0,0 @@
<!--Copyright 2024 Marigold authors and The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Marigold Pipelines for Computer Vision Tasks
![marigold](https://marigoldmonodepth.github.io/images/teaser_collage_compressed.jpg)
Marigold was proposed in [Repurposing Diffusion-Based Image Generators for Monocular Depth Estimation](https://huggingface.co/papers/2312.02145), a CVPR 2024 Oral paper by [Bingxin Ke](http://www.kebingxin.com/), [Anton Obukhov](https://www.obukhov.ai/), [Shengyu Huang](https://shengyuh.github.io/), [Nando Metzger](https://nandometzger.github.io/), [Rodrigo Caye Daudt](https://rcdaudt.github.io/), and [Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
The idea is to repurpose the rich generative prior of Text-to-Image Latent Diffusion Models (LDMs) for traditional computer vision tasks.
Initially, this idea was explored to fine-tune Stable Diffusion for Monocular Depth Estimation, as shown in the teaser above.
Later,
- [Tianfu Wang](https://tianfwang.github.io/) trained the first Latent Consistency Model (LCM) of Marigold, which unlocked fast single-step inference;
- [Kevin Qu](https://www.linkedin.com/in/kevin-qu-b3417621b/?locale=en_US) extended the approach to Surface Normals Estimation;
- [Anton Obukhov](https://www.obukhov.ai/) contributed the pipelines and documentation into diffusers (enabled and supported by [YiYi Xu](https://yiyixuxu.github.io/) and [Sayak Paul](https://sayak.dev/)).
The abstract from the paper is:
*Monocular depth estimation is a fundamental computer vision task. Recovering 3D depth from a single image is geometrically ill-posed and requires scene understanding, so it is not surprising that the rise of deep learning has led to a breakthrough. The impressive progress of monocular depth estimators has mirrored the growth in model capacity, from relatively modest CNNs to large Transformer architectures. Still, monocular depth estimators tend to struggle when presented with images with unfamiliar content and layout, since their knowledge of the visual world is restricted by the data seen during training, and challenged by zero-shot generalization to new domains. This motivates us to explore whether the extensive priors captured in recent generative diffusion models can enable better, more generalizable depth estimation. We introduce Marigold, a method for affine-invariant monocular depth estimation that is derived from Stable Diffusion and retains its rich prior knowledge. The estimator can be fine-tuned in a couple of days on a single GPU using only synthetic training data. It delivers state-of-the-art performance across a wide range of datasets, including over 20% performance gains in specific cases. Project page: https://marigoldmonodepth.github.io.*
## Available Pipelines
Each pipeline supports one Computer Vision task, which takes an input RGB image as input and produces a *prediction* of the modality of interest, such as a depth map of the input image.
Currently, the following tasks are implemented:
| Pipeline | Predicted Modalities | Demos |
|---------------------------------------------------------------------------------------------------------------------------------------------|------------------------------------------------------------------------------------------------------------------|:--------------------------------------------------------------------------------------------------------------------------------------------------:|
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-lcm), [Slow Original Demo (DDIM)](https://huggingface.co/spaces/prs-eth/marigold) |
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-normals-lcm) |
## Available Checkpoints
The original checkpoints can be found under the [PRS-ETH](https://huggingface.co/prs-eth/) Hugging Face organization.
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines. Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section [here](../../using-diffusers/svd#reduce-memory-usage).
</Tip>
<Tip warning={true}>
Marigold pipelines were designed and tested only with `DDIMScheduler` and `LCMScheduler`.
Depending on the scheduler, the number of inference steps required to get reliable predictions varies, and there is no universal value that works best across schedulers.
Because of that, the default value of `num_inference_steps` in the `__call__` method of the pipeline is set to `None` (see the API reference).
Unless set explicitly, its value will be taken from the checkpoint configuration `model_index.json`.
This is done to ensure high-quality predictions when calling the pipeline with just the `image` argument.
</Tip>
See also Marigold [usage examples](marigold_usage).
## MarigoldDepthPipeline
[[autodoc]] MarigoldDepthPipeline
- all
- __call__
## MarigoldNormalsPipeline
[[autodoc]] MarigoldNormalsPipeline
- all
- __call__
## MarigoldDepthOutput
[[autodoc]] pipelines.marigold.pipeline_marigold_depth.MarigoldDepthOutput
## MarigoldNormalsOutput
[[autodoc]] pipelines.marigold.pipeline_marigold_normals.MarigoldNormalsOutput

View File

@@ -57,7 +57,6 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Latent Consistency Models](latent_consistency_models) | text2image |
| [Latent Diffusion](latent_diffusion) | text2image, super-resolution |
| [LDM3D](stable_diffusion/ldm3d_diffusion) | text2image, text-to-3D, text-to-pano, upscaling |
| [LEDITS++](ledits_pp) | image editing |
| [MultiDiffusion](panorama) | text2image |
| [MusicLDM](musicldm) | text2audio |
| [Paint by Example](paint_by_example) | inpainting |
@@ -97,11 +96,6 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
- to
- components
[[autodoc]] pipelines.StableDiffusionMixin.enable_freeu
[[autodoc]] pipelines.StableDiffusionMixin.disable_freeu
## FlaxDiffusionPipeline
[[autodoc]] pipelines.pipeline_flax_utils.FlaxDiffusionPipeline

View File

@@ -31,13 +31,13 @@ Some notes about this pipeline:
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Inference with under 8GB GPU VRAM
Run the [`PixArtAlphaPipeline`] with under 8GB GPU VRAM by loading the text encoder in 8-bit precision. Let's walk through a full-fledged example.
Run the [`PixArtAlphaPipeline`] with under 8GB GPU VRAM by loading the text encoder in 8-bit precision. Let's walk through a full-fledged example.
First, install the [bitsandbytes](https://github.com/TimDettmers/bitsandbytes) library:
@@ -75,10 +75,10 @@ with torch.no_grad():
prompt_embeds, prompt_attention_mask, negative_embeds, negative_prompt_attention_mask = pipe.encode_prompt(prompt)
```
Since text embeddings have been computed, remove the `text_encoder` and `pipe` from the memory, and free up some GPU VRAM:
Since text embeddings have been computed, remove the `text_encoder` and `pipe` from the memory, and free up som GPU VRAM:
```python
import gc
import gc
def flush():
gc.collect()
@@ -99,7 +99,7 @@ pipe = PixArtAlphaPipeline.from_pretrained(
).to("cuda")
latents = pipe(
negative_prompt=None,
negative_prompt=None,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_embeds,
prompt_attention_mask=prompt_attention_mask,
@@ -146,3 +146,4 @@ While loading the `text_encoder`, you set `load_in_8bit` to `True`. You could al
[[autodoc]] PixArtAlphaPipeline
- all
- __call__

View File

@@ -1,149 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# PixArt-Σ
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/pixart/header_collage_sigma.jpg)
[PixArt-Σ: Weak-to-Strong Training of Diffusion Transformer for 4K Text-to-Image Generation](https://huggingface.co/papers/2403.04692) is Junsong Chen, Jincheng Yu, Chongjian Ge, Lewei Yao, Enze Xie, Yue Wu, Zhongdao Wang, James Kwok, Ping Luo, Huchuan Lu, and Zhenguo Li.
The abstract from the paper is:
*In this paper, we introduce PixArt-Σ, a Diffusion Transformer model (DiT) capable of directly generating images at 4K resolution. PixArt-Σ represents a significant advancement over its predecessor, PixArt-α, offering images of markedly higher fidelity and improved alignment with text prompts. A key feature of PixArt-Σ is its training efficiency. Leveraging the foundational pre-training of PixArt-α, it evolves from the weaker baseline to a stronger model via incorporating higher quality data, a process we term “weak-to-strong training”. The advancements in PixArt-Σ are twofold: (1) High-Quality Training Data: PixArt-Σ incorporates superior-quality image data, paired with more precise and detailed image captions. (2) Efficient Token Compression: we propose a novel attention module within the DiT framework that compresses both keys and values, significantly improving efficiency and facilitating ultra-high-resolution image generation. Thanks to these improvements, PixArt-Σ achieves superior image quality and user prompt adherence capabilities with significantly smaller model size (0.6B parameters) than existing text-to-image diffusion models, such as SDXL (2.6B parameters) and SD Cascade (5.1B parameters). Moreover, PixArt-Σs capability to generate 4K images supports the creation of high-resolution posters and wallpapers, efficiently bolstering the production of highquality visual content in industries such as film and gaming.*
You can find the original codebase at [PixArt-alpha/PixArt-sigma](https://github.com/PixArt-alpha/PixArt-sigma) and all the available checkpoints at [PixArt-alpha](https://huggingface.co/PixArt-alpha).
Some notes about this pipeline:
* It uses a Transformer backbone (instead of a UNet) for denoising. As such it has a similar architecture as [DiT](https://hf.co/docs/transformers/model_doc/dit).
* It was trained using text conditions computed from T5. This aspect makes the pipeline better at following complex text prompts with intricate details.
* It is good at producing high-resolution images at different aspect ratios. To get the best results, the authors recommend some size brackets which can be found [here](https://github.com/PixArt-alpha/PixArt-sigma/blob/master/diffusion/data/datasets/utils.py).
* It rivals the quality of state-of-the-art text-to-image generation systems (as of this writing) such as PixArt-α, Stable Diffusion XL, Playground V2.0 and DALL-E 3, while being more efficient than them.
* It shows the ability of generating super high resolution images, such as 2048px or even 4K.
* It shows that text-to-image models can grow from a weak model to a stronger one through several improvements (VAEs, datasets, and so on.)
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Inference with under 8GB GPU VRAM
Run the [`PixArtSigmaPipeline`] with under 8GB GPU VRAM by loading the text encoder in 8-bit precision. Let's walk through a full-fledged example.
First, install the [bitsandbytes](https://github.com/TimDettmers/bitsandbytes) library:
```bash
pip install -U bitsandbytes
```
Then load the text encoder in 8-bit:
```python
from transformers import T5EncoderModel
from diffusers import PixArtSigmaPipeline
import torch
text_encoder = T5EncoderModel.from_pretrained(
"PixArt-alpha/PixArt-Sigma-XL-2-1024-MS",
subfolder="text_encoder",
load_in_8bit=True,
device_map="auto",
)
pipe = PixArtSigmaPipeline.from_pretrained(
"PixArt-alpha/PixArt-Sigma-XL-2-1024-MS",
text_encoder=text_encoder,
transformer=None,
device_map="balanced"
)
```
Now, use the `pipe` to encode a prompt:
```python
with torch.no_grad():
prompt = "cute cat"
prompt_embeds, prompt_attention_mask, negative_embeds, negative_prompt_attention_mask = pipe.encode_prompt(prompt)
```
Since text embeddings have been computed, remove the `text_encoder` and `pipe` from the memory, and free up some GPU VRAM:
```python
import gc
def flush():
gc.collect()
torch.cuda.empty_cache()
del text_encoder
del pipe
flush()
```
Then compute the latents with the prompt embeddings as inputs:
```python
pipe = PixArtSigmaPipeline.from_pretrained(
"PixArt-alpha/PixArt-Sigma-XL-2-1024-MS",
text_encoder=None,
torch_dtype=torch.float16,
).to("cuda")
latents = pipe(
negative_prompt=None,
prompt_embeds=prompt_embeds,
negative_prompt_embeds=negative_embeds,
prompt_attention_mask=prompt_attention_mask,
negative_prompt_attention_mask=negative_prompt_attention_mask,
num_images_per_prompt=1,
output_type="latent",
).images
del pipe.transformer
flush()
```
<Tip>
Notice that while initializing `pipe`, you're setting `text_encoder` to `None` so that it's not loaded.
</Tip>
Once the latents are computed, pass it off to the VAE to decode into a real image:
```python
with torch.no_grad():
image = pipe.vae.decode(latents / pipe.vae.config.scaling_factor, return_dict=False)[0]
image = pipe.image_processor.postprocess(image, output_type="pil")[0]
image.save("cat.png")
```
By deleting components you aren't using and flushing the GPU VRAM, you should be able to run [`PixArtSigmaPipeline`] with under 8GB GPU VRAM.
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/pixart/8bits_cat.png)
If you want a report of your memory-usage, run this [script](https://gist.github.com/sayakpaul/3ae0f847001d342af27018a96f467e4e).
<Tip warning={true}>
Text embeddings computed in 8-bit can impact the quality of the generated images because of the information loss in the representation space caused by the reduced precision. It's recommended to compare the outputs with and without 8-bit.
</Tip>
While loading the `text_encoder`, you set `load_in_8bit` to `True`. You could also specify `load_in_4bit` to bring your memory requirements down even further to under 7GB.
## PixArtSigmaPipeline
[[autodoc]] PixArtSigmaPipeline
- all
- __call__

View File

@@ -30,6 +30,6 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
- all
- __call__
## SemanticStableDiffusionPipelineOutput
## StableDiffusionSafePipelineOutput
[[autodoc]] pipelines.semantic_stable_diffusion.pipeline_output.SemanticStableDiffusionPipelineOutput
- all

View File

@@ -12,13 +12,13 @@ specific language governing permissions and limitations under the License.
# Stable Cascade
This model is built upon the [Würstchen](https://openreview.net/forum?id=gU58d5QeGv) architecture and its main
difference to other models like Stable Diffusion is that it is working at a much smaller latent space. Why is this
important? The smaller the latent space, the **faster** you can run inference and the **cheaper** the training becomes.
How small is the latent space? Stable Diffusion uses a compression factor of 8, resulting in a 1024x1024 image being
encoded to 128x128. Stable Cascade achieves a compression factor of 42, meaning that it is possible to encode a
1024x1024 image to 24x24, while maintaining crisp reconstructions. The text-conditional model is then trained in the
highly compressed latent space. Previous versions of this architecture, achieved a 16x cost reduction over Stable
This model is built upon the [Würstchen](https://openreview.net/forum?id=gU58d5QeGv) architecture and its main
difference to other models like Stable Diffusion is that it is working at a much smaller latent space. Why is this
important? The smaller the latent space, the **faster** you can run inference and the **cheaper** the training becomes.
How small is the latent space? Stable Diffusion uses a compression factor of 8, resulting in a 1024x1024 image being
encoded to 128x128. Stable Cascade achieves a compression factor of 42, meaning that it is possible to encode a
1024x1024 image to 24x24, while maintaining crisp reconstructions. The text-conditional model is then trained in the
highly compressed latent space. Previous versions of this architecture, achieved a 16x cost reduction over Stable
Diffusion 1.5.
Therefore, this kind of model is well suited for usages where efficiency is important. Furthermore, all known extensions
@@ -30,154 +30,13 @@ The original codebase can be found at [Stability-AI/StableCascade](https://githu
Stable Cascade consists of three models: Stage A, Stage B and Stage C, representing a cascade to generate images,
hence the name "Stable Cascade".
Stage A & B are used to compress images, similar to what the job of the VAE is in Stable Diffusion.
However, with this setup, a much higher compression of images can be achieved. While the Stable Diffusion models use a
spatial compression factor of 8, encoding an image with resolution of 1024 x 1024 to 128 x 128, Stable Cascade achieves
a compression factor of 42. This encodes a 1024 x 1024 image to 24 x 24, while being able to accurately decode the
image. This comes with the great benefit of cheaper training and inference. Furthermore, Stage C is responsible
Stage A & B are used to compress images, similar to what the job of the VAE is in Stable Diffusion.
However, with this setup, a much higher compression of images can be achieved. While the Stable Diffusion models use a
spatial compression factor of 8, encoding an image with resolution of 1024 x 1024 to 128 x 128, Stable Cascade achieves
a compression factor of 42. This encodes a 1024 x 1024 image to 24 x 24, while being able to accurately decode the
image. This comes with the great benefit of cheaper training and inference. Furthermore, Stage C is responsible
for generating the small 24 x 24 latents given a text prompt.
The Stage C model operates on the small 24 x 24 latents and denoises the latents conditioned on text prompts. The model is also the largest component in the Cascade pipeline and is meant to be used with the `StableCascadePriorPipeline`
The Stage B and Stage A models are used with the `StableCascadeDecoderPipeline` and are responsible for generating the final image given the small 24 x 24 latents.
<Tip warning={true}>
There are some restrictions on data types that can be used with the Stable Cascade models. The official checkpoints for the `StableCascadePriorPipeline` do not support the `torch.float16` data type. Please use `torch.bfloat16` instead.
In order to use the `torch.bfloat16` data type with the `StableCascadeDecoderPipeline` you need to have PyTorch 2.2.0 or higher installed. This also means that using the `StableCascadeCombinedPipeline` with `torch.bfloat16` requires PyTorch 2.2.0 or higher, since it calls the `StableCascadeDecoderPipeline` internally.
If it is not possible to install PyTorch 2.2.0 or higher in your environment, the `StableCascadeDecoderPipeline` can be used on its own with the `torch.float16` data type. You can download the full precision or `bf16` variant weights for the pipeline and cast the weights to `torch.float16`.
</Tip>
## Usage example
```python
import torch
from diffusers import StableCascadeDecoderPipeline, StableCascadePriorPipeline
prompt = "an image of a shiba inu, donning a spacesuit and helmet"
negative_prompt = ""
prior = StableCascadePriorPipeline.from_pretrained("stabilityai/stable-cascade-prior", variant="bf16", torch_dtype=torch.bfloat16)
decoder = StableCascadeDecoderPipeline.from_pretrained("stabilityai/stable-cascade", variant="bf16", torch_dtype=torch.float16)
prior.enable_model_cpu_offload()
prior_output = prior(
prompt=prompt,
height=1024,
width=1024,
negative_prompt=negative_prompt,
guidance_scale=4.0,
num_images_per_prompt=1,
num_inference_steps=20
)
decoder.enable_model_cpu_offload()
decoder_output = decoder(
image_embeddings=prior_output.image_embeddings.to(torch.float16),
prompt=prompt,
negative_prompt=negative_prompt,
guidance_scale=0.0,
output_type="pil",
num_inference_steps=10
).images[0]
decoder_output.save("cascade.png")
```
## Using the Lite Versions of the Stage B and Stage C models
```python
import torch
from diffusers import (
StableCascadeDecoderPipeline,
StableCascadePriorPipeline,
StableCascadeUNet,
)
prompt = "an image of a shiba inu, donning a spacesuit and helmet"
negative_prompt = ""
prior_unet = StableCascadeUNet.from_pretrained("stabilityai/stable-cascade-prior", subfolder="prior_lite")
decoder_unet = StableCascadeUNet.from_pretrained("stabilityai/stable-cascade", subfolder="decoder_lite")
prior = StableCascadePriorPipeline.from_pretrained("stabilityai/stable-cascade-prior", prior=prior_unet)
decoder = StableCascadeDecoderPipeline.from_pretrained("stabilityai/stable-cascade", decoder=decoder_unet)
prior.enable_model_cpu_offload()
prior_output = prior(
prompt=prompt,
height=1024,
width=1024,
negative_prompt=negative_prompt,
guidance_scale=4.0,
num_images_per_prompt=1,
num_inference_steps=20
)
decoder.enable_model_cpu_offload()
decoder_output = decoder(
image_embeddings=prior_output.image_embeddings,
prompt=prompt,
negative_prompt=negative_prompt,
guidance_scale=0.0,
output_type="pil",
num_inference_steps=10
).images[0]
decoder_output.save("cascade.png")
```
## Loading original checkpoints with `from_single_file`
Loading the original format checkpoints is supported via `from_single_file` method in the StableCascadeUNet.
```python
import torch
from diffusers import (
StableCascadeDecoderPipeline,
StableCascadePriorPipeline,
StableCascadeUNet,
)
prompt = "an image of a shiba inu, donning a spacesuit and helmet"
negative_prompt = ""
prior_unet = StableCascadeUNet.from_single_file(
"https://huggingface.co/stabilityai/stable-cascade/resolve/main/stage_c_bf16.safetensors",
torch_dtype=torch.bfloat16
)
decoder_unet = StableCascadeUNet.from_single_file(
"https://huggingface.co/stabilityai/stable-cascade/blob/main/stage_b_bf16.safetensors",
torch_dtype=torch.bfloat16
)
prior = StableCascadePriorPipeline.from_pretrained("stabilityai/stable-cascade-prior", prior=prior_unet, torch_dtype=torch.bfloat16)
decoder = StableCascadeDecoderPipeline.from_pretrained("stabilityai/stable-cascade", decoder=decoder_unet, torch_dtype=torch.bfloat16)
prior.enable_model_cpu_offload()
prior_output = prior(
prompt=prompt,
height=1024,
width=1024,
negative_prompt=negative_prompt,
guidance_scale=4.0,
num_images_per_prompt=1,
num_inference_steps=20
)
decoder.enable_model_cpu_offload()
decoder_output = decoder(
image_embeddings=prior_output.image_embeddings,
prompt=prompt,
negative_prompt=negative_prompt,
guidance_scale=0.0,
output_type="pil",
num_inference_steps=10
).images[0]
decoder_output.save("cascade-single-file.png")
```
## Uses
### Direct Use
@@ -194,7 +53,7 @@ Excluded uses are described below.
### Out-of-Scope Use
The model was not trained to be factual or true representations of people or events,
The model was not trained to be factual or true representations of people or events,
and therefore using the model to generate such content is out-of-scope for the abilities of this model.
The model should not be used in any way that violates Stability AI's [Acceptable Use Policy](https://stability.ai/use-policy).

View File

@@ -10,7 +10,9 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# T2I-Adapter
# Text-to-Image Generation with Adapter Conditioning
## Overview
[T2I-Adapter: Learning Adapters to Dig out More Controllable Ability for Text-to-Image Diffusion Models](https://arxiv.org/abs/2302.08453) by Chong Mou, Xintao Wang, Liangbin Xie, Jian Zhang, Zhongang Qi, Ying Shan, Xiaohu Qie.
@@ -22,26 +24,236 @@ The abstract of the paper is the following:
This model was contributed by the community contributor [HimariO](https://github.com/HimariO) ❤️ .
## StableDiffusionAdapterPipeline
## Available Pipelines:
| Pipeline | Tasks | Demo
|---|---|:---:|
| [StableDiffusionAdapterPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_adapter.py) | *Text-to-Image Generation with T2I-Adapter Conditioning* | -
| [StableDiffusionXLAdapterPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_xl_adapter.py) | *Text-to-Image Generation with T2I-Adapter Conditioning on StableDiffusion-XL* | -
## Usage example with the base model of StableDiffusion-1.4/1.5
In the following we give a simple example of how to use a *T2I-Adapter* checkpoint with Diffusers for inference based on StableDiffusion-1.4/1.5.
All adapters use the same pipeline.
1. Images are first converted into the appropriate *control image* format.
2. The *control image* and *prompt* are passed to the [`StableDiffusionAdapterPipeline`].
Let's have a look at a simple example using the [Color Adapter](https://huggingface.co/TencentARC/t2iadapter_color_sd14v1).
```python
from diffusers.utils import load_image, make_image_grid
image = load_image("https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_ref.png")
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_ref.png)
Then we can create our color palette by simply resizing it to 8 by 8 pixels and then scaling it back to original size.
```python
from PIL import Image
color_palette = image.resize((8, 8))
color_palette = color_palette.resize((512, 512), resample=Image.Resampling.NEAREST)
```
Let's take a look at the processed image.
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_palette.png)
Next, create the adapter pipeline
```py
import torch
from diffusers import StableDiffusionAdapterPipeline, T2IAdapter
adapter = T2IAdapter.from_pretrained("TencentARC/t2iadapter_color_sd14v1", torch_dtype=torch.float16)
pipe = StableDiffusionAdapterPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
adapter=adapter,
torch_dtype=torch.float16,
)
pipe.to("cuda")
```
Finally, pass the prompt and control image to the pipeline
```py
# fix the random seed, so you will get the same result as the example
generator = torch.Generator("cuda").manual_seed(7)
out_image = pipe(
"At night, glowing cubes in front of the beach",
image=color_palette,
generator=generator,
).images[0]
make_image_grid([image, color_palette, out_image], rows=1, cols=3)
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_output.png)
## Usage example with the base model of StableDiffusion-XL
In the following we give a simple example of how to use a *T2I-Adapter* checkpoint with Diffusers for inference based on StableDiffusion-XL.
All adapters use the same pipeline.
1. Images are first downloaded into the appropriate *control image* format.
2. The *control image* and *prompt* are passed to the [`StableDiffusionXLAdapterPipeline`].
Let's have a look at a simple example using the [Sketch Adapter](https://huggingface.co/Adapter/t2iadapter/tree/main/sketch_sdxl_1.0).
```python
from diffusers.utils import load_image, make_image_grid
sketch_image = load_image("https://huggingface.co/Adapter/t2iadapter/resolve/main/sketch.png").convert("L")
```
![img](https://huggingface.co/Adapter/t2iadapter/resolve/main/sketch.png)
Then, create the adapter pipeline
```py
import torch
from diffusers import (
T2IAdapter,
StableDiffusionXLAdapterPipeline,
DDPMScheduler
)
model_id = "stabilityai/stable-diffusion-xl-base-1.0"
adapter = T2IAdapter.from_pretrained("Adapter/t2iadapter", subfolder="sketch_sdxl_1.0", torch_dtype=torch.float16, adapter_type="full_adapter_xl")
scheduler = DDPMScheduler.from_pretrained(model_id, subfolder="scheduler")
pipe = StableDiffusionXLAdapterPipeline.from_pretrained(
model_id, adapter=adapter, safety_checker=None, torch_dtype=torch.float16, variant="fp16", scheduler=scheduler
)
pipe.to("cuda")
```
Finally, pass the prompt and control image to the pipeline
```py
# fix the random seed, so you will get the same result as the example
generator = torch.Generator().manual_seed(42)
sketch_image_out = pipe(
prompt="a photo of a dog in real world, high quality",
negative_prompt="extra digit, fewer digits, cropped, worst quality, low quality",
image=sketch_image,
generator=generator,
guidance_scale=7.5
).images[0]
make_image_grid([sketch_image, sketch_image_out], rows=1, cols=2)
```
![img](https://huggingface.co/Adapter/t2iadapter/resolve/main/sketch_output.png)
## Available checkpoints
Non-diffusers checkpoints can be found under [TencentARC/T2I-Adapter](https://huggingface.co/TencentARC/T2I-Adapter/tree/main/models).
### T2I-Adapter with Stable Diffusion 1.4
| Model Name | Control Image Overview| Control Image Example | Generated Image Example |
|---|---|---|---|
|[TencentARC/t2iadapter_color_sd14v1](https://huggingface.co/TencentARC/t2iadapter_color_sd14v1)<br/> *Trained with spatial color palette* | An image with 8x8 color palette.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_sample_input.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_sample_output.png"/></a>|
|[TencentARC/t2iadapter_canny_sd14v1](https://huggingface.co/TencentARC/t2iadapter_canny_sd14v1)<br/> *Trained with canny edge detection* | A monochrome image with white edges on a black background.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/canny_sample_input.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/canny_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/canny_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/canny_sample_output.png"/></a>|
|[TencentARC/t2iadapter_sketch_sd14v1](https://huggingface.co/TencentARC/t2iadapter_sketch_sd14v1)<br/> *Trained with [PidiNet](https://github.com/zhuoinoulu/pidinet) edge detection* | A hand-drawn monochrome image with white outlines on a black background.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/sketch_sample_input.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/sketch_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/sketch_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/sketch_sample_output.png"/></a>|
|[TencentARC/t2iadapter_depth_sd14v1](https://huggingface.co/TencentARC/t2iadapter_depth_sd14v1)<br/> *Trained with Midas depth estimation* | A grayscale image with black representing deep areas and white representing shallow areas.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_input.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_output.png"/></a>|
|[TencentARC/t2iadapter_openpose_sd14v1](https://huggingface.co/TencentARC/t2iadapter_openpose_sd14v1)<br/> *Trained with OpenPose bone image* | A [OpenPose bone](https://github.com/CMU-Perceptual-Computing-Lab/openpose) image.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/openpose_sample_input.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/openpose_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/openpose_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/openpose_sample_output.png"/></a>|
|[TencentARC/t2iadapter_keypose_sd14v1](https://huggingface.co/TencentARC/t2iadapter_keypose_sd14v1)<br/> *Trained with mmpose skeleton image* | A [mmpose skeleton](https://github.com/open-mmlab/mmpose) image.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_input.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_output.png"/></a>|
|[TencentARC/t2iadapter_seg_sd14v1](https://huggingface.co/TencentARC/t2iadapter_seg_sd14v1)<br/>*Trained with semantic segmentation* | An [custom](https://github.com/TencentARC/T2I-Adapter/discussions/25) segmentation protocol image.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/seg_sample_input.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/seg_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/seg_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/seg_sample_output.png"/></a> |
|[TencentARC/t2iadapter_canny_sd15v2](https://huggingface.co/TencentARC/t2iadapter_canny_sd15v2)||
|[TencentARC/t2iadapter_depth_sd15v2](https://huggingface.co/TencentARC/t2iadapter_depth_sd15v2)||
|[TencentARC/t2iadapter_sketch_sd15v2](https://huggingface.co/TencentARC/t2iadapter_sketch_sd15v2)||
|[TencentARC/t2iadapter_zoedepth_sd15v1](https://huggingface.co/TencentARC/t2iadapter_zoedepth_sd15v1)||
|[Adapter/t2iadapter, subfolder='sketch_sdxl_1.0'](https://huggingface.co/Adapter/t2iadapter/tree/main/sketch_sdxl_1.0)||
|[Adapter/t2iadapter, subfolder='canny_sdxl_1.0'](https://huggingface.co/Adapter/t2iadapter/tree/main/canny_sdxl_1.0)||
|[Adapter/t2iadapter, subfolder='openpose_sdxl_1.0'](https://huggingface.co/Adapter/t2iadapter/tree/main/openpose_sdxl_1.0)||
## Combining multiple adapters
[`MultiAdapter`] can be used for applying multiple conditionings at once.
Here we use the keypose adapter for the character posture and the depth adapter for creating the scene.
```py
from diffusers.utils import load_image, make_image_grid
cond_keypose = load_image(
"https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_input.png"
)
cond_depth = load_image(
"https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_input.png"
)
cond = [cond_keypose, cond_depth]
prompt = ["A man walking in an office room with a nice view"]
```
The two control images look as such:
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_input.png)
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_input.png)
`MultiAdapter` combines keypose and depth adapters.
`adapter_conditioning_scale` balances the relative influence of the different adapters.
```py
import torch
from diffusers import StableDiffusionAdapterPipeline, MultiAdapter, T2IAdapter
adapters = MultiAdapter(
[
T2IAdapter.from_pretrained("TencentARC/t2iadapter_keypose_sd14v1"),
T2IAdapter.from_pretrained("TencentARC/t2iadapter_depth_sd14v1"),
]
)
adapters = adapters.to(torch.float16)
pipe = StableDiffusionAdapterPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
torch_dtype=torch.float16,
adapter=adapters,
).to("cuda")
image = pipe(prompt, cond, adapter_conditioning_scale=[0.8, 0.8]).images[0]
make_image_grid([cond_keypose, cond_depth, image], rows=1, cols=3)
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_depth_sample_output.png)
## T2I-Adapter vs ControlNet
T2I-Adapter is similar to [ControlNet](https://huggingface.co/docs/diffusers/main/en/api/pipelines/controlnet).
T2I-Adapter uses a smaller auxiliary network which is only run once for the entire diffusion process.
However, T2I-Adapter performs slightly worse than ControlNet.
## StableDiffusionAdapterPipeline
[[autodoc]] StableDiffusionAdapterPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_vae_slicing
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_vae_slicing
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
## StableDiffusionXLAdapterPipeline
[[autodoc]] StableDiffusionXLAdapterPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_vae_slicing
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_vae_slicing
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention

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@@ -172,41 +172,3 @@ inpaint = StableDiffusionInpaintPipeline(**text2img.components)
# now you can use text2img(...), img2img(...), inpaint(...) just like the call methods of each respective pipeline
```
### Create web demos using `gradio`
The Stable Diffusion pipelines are automatically supported in [Gradio](https://github.com/gradio-app/gradio/), a library that makes creating beautiful and user-friendly machine learning apps on the web a breeze. First, make sure you have Gradio installed:
```sh
pip install -U gradio
```
Then, create a web demo around any Stable Diffusion-based pipeline. For example, you can create an image generation pipeline in a single line of code with Gradio's [`Interface.from_pipeline`](https://www.gradio.app/docs/interface#interface-from-pipeline) function:
```py
from diffusers import StableDiffusionPipeline
import gradio as gr
pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4")
gr.Interface.from_pipeline(pipe).launch()
```
which opens an intuitive drag-and-drop interface in your browser:
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/gradio-panda.png)
Similarly, you could create a demo for an image-to-image pipeline with:
```py
from diffusers import StableDiffusionImg2ImgPipeline
import gradio as gr
pipe = StableDiffusionImg2ImgPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
gr.Interface.from_pipeline(pipe).launch()
```
By default, the web demo runs on a local server. If you'd like to share it with others, you can generate a temporary public
link by setting `share=True` in `launch()`. Or, you can host your demo on [Hugging Face Spaces](https://huggingface.co/spaces)https://huggingface.co/spaces for a permanent link.

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@@ -1,251 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Stable Diffusion 3
Stable Diffusion 3 (SD3) was proposed in [Scaling Rectified Flow Transformers for High-Resolution Image Synthesis](https://arxiv.org/pdf/2403.03206.pdf) by Patrick Esser, Sumith Kulal, Andreas Blattmann, Rahim Entezari, Jonas Muller, Harry Saini, Yam Levi, Dominik Lorenz, Axel Sauer, Frederic Boesel, Dustin Podell, Tim Dockhorn, Zion English, Kyle Lacey, Alex Goodwin, Yannik Marek, and Robin Rombach.
The abstract from the paper is:
*Diffusion models create data from noise by inverting the forward paths of data towards noise and have emerged as a powerful generative modeling technique for high-dimensional, perceptual data such as images and videos. Rectified flow is a recent generative model formulation that connects data and noise in a straight line. Despite its better theoretical properties and conceptual simplicity, it is not yet decisively established as standard practice. In this work, we improve existing noise sampling techniques for training rectified flow models by biasing them towards perceptually relevant scales. Through a large-scale study, we demonstrate the superior performance of this approach compared to established diffusion formulations for high-resolution text-to-image synthesis. Additionally, we present a novel transformer-based architecture for text-to-image generation that uses separate weights for the two modalities and enables a bidirectional flow of information between image and text tokens, improving text comprehension typography, and human preference ratings. We demonstrate that this architecture follows predictable scaling trends and correlates lower validation loss to improved text-to-image synthesis as measured by various metrics and human evaluations.*
## Usage Example
_As the model is gated, before using it with diffusers you first need to go to the [Stable Diffusion 3 Medium Hugging Face page](https://huggingface.co/stabilityai/stable-diffusion-3-medium-diffusers), fill in the form and accept the gate. Once you are in, you need to login so that your system knows youve accepted the gate._
Use the command below to log in:
```bash
huggingface-cli login
```
<Tip>
The SD3 pipeline uses three text encoders to generate an image. Model offloading is necessary in order for it to run on most commodity hardware. Please use the `torch.float16` data type for additional memory savings.
</Tip>
```python
import torch
from diffusers import StableDiffusion3Pipeline
pipe = StableDiffusion3Pipeline.from_pretrained("stabilityai/stable-diffusion-3-medium-diffusers", torch_dtype=torch.float16)
pipe.to("cuda")
image = pipe(
prompt="a photo of a cat holding a sign that says hello world",
negative_prompt="",
num_inference_steps=28,
height=1024,
width=1024,
guidance_scale=7.0,
).images[0]
image.save("sd3_hello_world.png")
```
## Memory Optimisations for SD3
SD3 uses three text encoders, one if which is the very large T5-XXL model. This makes it challenging to run the model on GPUs with less than 24GB of VRAM, even when using `fp16` precision. The following section outlines a few memory optimizations in Diffusers that make it easier to run SD3 on low resource hardware.
### Running Inference with Model Offloading
The most basic memory optimization available in Diffusers allows you to offload the components of the model to CPU during inference in order to save memory, while seeing a slight increase in inference latency. Model offloading will only move a model component onto the GPU when it needs to be executed, while keeping the remaining components on the CPU.
```python
import torch
from diffusers import StableDiffusion3Pipeline
pipe = StableDiffusion3Pipeline.from_pretrained("stabilityai/stable-diffusion-3-medium-diffusers", torch_dtype=torch.float16)
pipe.enable_model_cpu_offload()
image = pipe(
prompt="a photo of a cat holding a sign that says hello world",
negative_prompt="",
num_inference_steps=28,
height=1024,
width=1024,
guidance_scale=7.0,
).images[0]
image.save("sd3_hello_world.png")
```
### Dropping the T5 Text Encoder during Inference
Removing the memory-intensive 4.7B parameter T5-XXL text encoder during inference can significantly decrease the memory requirements for SD3 with only a slight loss in performance.
```python
import torch
from diffusers import StableDiffusion3Pipeline
pipe = StableDiffusion3Pipeline.from_pretrained(
"stabilityai/stable-diffusion-3-medium-diffusers",
text_encoder_3=None,
tokenizer_3=None,
torch_dtype=torch.float16
)
pipe.to("cuda")
image = pipe(
prompt="a photo of a cat holding a sign that says hello world",
negative_prompt="",
num_inference_steps=28,
height=1024,
width=1024,
guidance_scale=7.0,
).images[0]
image.save("sd3_hello_world-no-T5.png")
```
### Using a Quantized Version of the T5 Text Encoder
We can leverage the `bitsandbytes` library to load and quantize the T5-XXL text encoder to 8-bit precision. This allows you to keep using all three text encoders while only slightly impacting performance.
First install the `bitsandbytes` library.
```shell
pip install bitsandbytes
```
Then load the T5-XXL model using the `BitsAndBytesConfig`.
```python
import torch
from diffusers import StableDiffusion3Pipeline
from transformers import T5EncoderModel, BitsAndBytesConfig
quantization_config = BitsAndBytesConfig(load_in_8bit=True)
model_id = "stabilityai/stable-diffusion-3-medium-diffusers"
text_encoder = T5EncoderModel.from_pretrained(
model_id,
subfolder="text_encoder_3",
quantization_config=quantization_config,
)
pipe = StableDiffusion3Pipeline.from_pretrained(
model_id,
text_encoder_3=text_encoder,
device_map="balanced",
torch_dtype=torch.float16
)
image = pipe(
prompt="a photo of a cat holding a sign that says hello world",
negative_prompt="",
num_inference_steps=28,
height=1024,
width=1024,
guidance_scale=7.0,
).images[0]
image.save("sd3_hello_world-8bit-T5.png")
```
You can find the end-to-end script [here](https://gist.github.com/sayakpaul/82acb5976509851f2db1a83456e504f1).
## Performance Optimizations for SD3
### Using Torch Compile to Speed Up Inference
Using compiled components in the SD3 pipeline can speed up inference by as much as 4X. The following code snippet demonstrates how to compile the Transformer and VAE components of the SD3 pipeline.
```python
import torch
from diffusers import StableDiffusion3Pipeline
torch.set_float32_matmul_precision("high")
torch._inductor.config.conv_1x1_as_mm = True
torch._inductor.config.coordinate_descent_tuning = True
torch._inductor.config.epilogue_fusion = False
torch._inductor.config.coordinate_descent_check_all_directions = True
pipe = StableDiffusion3Pipeline.from_pretrained(
"stabilityai/stable-diffusion-3-medium-diffusers",
torch_dtype=torch.float16
).to("cuda")
pipe.set_progress_bar_config(disable=True)
pipe.transformer.to(memory_format=torch.channels_last)
pipe.vae.to(memory_format=torch.channels_last)
pipe.transformer = torch.compile(pipe.transformer, mode="max-autotune", fullgraph=True)
pipe.vae.decode = torch.compile(pipe.vae.decode, mode="max-autotune", fullgraph=True)
# Warm Up
prompt = "a photo of a cat holding a sign that says hello world"
for _ in range(3):
_ = pipe(prompt=prompt, generator=torch.manual_seed(1))
# Run Inference
image = pipe(prompt=prompt, generator=torch.manual_seed(1)).images[0]
image.save("sd3_hello_world.png")
```
Check out the full script [here](https://gist.github.com/sayakpaul/508d89d7aad4f454900813da5d42ca97).
## Loading the original checkpoints via `from_single_file`
The `SD3Transformer2DModel` and `StableDiffusion3Pipeline` classes support loading the original checkpoints via the `from_single_file` method. This method allows you to load the original checkpoint files that were used to train the models.
## Loading the original checkpoints for the `SD3Transformer2DModel`
```python
from diffusers import SD3Transformer2DModel
model = SD3Transformer2DModel.from_single_file("https://huggingface.co/stabilityai/stable-diffusion-3-medium/blob/main/sd3_medium.safetensors")
```
## Loading the single checkpoint for the `StableDiffusion3Pipeline`
### Loading the single file checkpoint without T5
```python
import torch
from diffusers import StableDiffusion3Pipeline
pipe = StableDiffusion3Pipeline.from_single_file(
"https://huggingface.co/stabilityai/stable-diffusion-3-medium/blob/main/sd3_medium_incl_clips.safetensors",
torch_dtype=torch.float16,
text_encoder_3=None
)
pipe.enable_model_cpu_offload()
image = pipe("a picture of a cat holding a sign that says hello world").images[0]
image.save('sd3-single-file.png')
```
### Loading the single file checkpoint without T5
```python
import torch
from diffusers import StableDiffusion3Pipeline
pipe = StableDiffusion3Pipeline.from_single_file(
"https://huggingface.co/stabilityai/stable-diffusion-3-medium/blob/main/sd3_medium_incl_clips_t5xxlfp8.safetensors",
torch_dtype=torch.float16,
)
pipe.enable_model_cpu_offload()
image = pipe("a picture of a cat holding a sign that says hello world").images[0]
image.save('sd3-single-file-t5-fp8.png')
```
## StableDiffusion3Pipeline
[[autodoc]] StableDiffusion3Pipeline
- all
- __call__

View File

@@ -1,22 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# EDMEulerScheduler
The Karras formulation of the Euler scheduler (Algorithm 2) from the [Elucidating the Design Space of Diffusion-Based Generative Models](https://huggingface.co/papers/2206.00364) paper by Karras et al. This is a fast scheduler which can often generate good outputs in 20-30 steps. The scheduler is based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L51) implementation by [Katherine Crowson](https://github.com/crowsonkb/).
## EDMEulerScheduler
[[autodoc]] EDMEulerScheduler
## EDMEulerSchedulerOutput
[[autodoc]] schedulers.scheduling_edm_euler.EDMEulerSchedulerOutput

View File

@@ -1,24 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# EDMDPMSolverMultistepScheduler
`EDMDPMSolverMultistepScheduler` is a [Karras formulation](https://huggingface.co/papers/2206.00364) of `DPMSolverMultistepScheduler`, a multistep scheduler from [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095) by Cheng Lu, Yuhao Zhou, Fan Bao, Jianfei Chen, Chongxuan Li, and Jun Zhu.
DPMSolver (and the improved version DPMSolver++) is a fast dedicated high-order solver for diffusion ODEs with convergence order guarantee. Empirically, DPMSolver sampling with only 20 steps can generate high-quality
samples, and it can generate quite good samples even in 10 steps.
## EDMDPMSolverMultistepScheduler
[[autodoc]] EDMDPMSolverMultistepScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput

View File

@@ -1,18 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# FlowMatchEulerDiscreteScheduler
`FlowMatchEulerDiscreteScheduler` is based on the flow-matching sampling introduced in [Stable Diffusion 3](https://arxiv.org/abs/2403.03206).
## FlowMatchEulerDiscreteScheduler
[[autodoc]] FlowMatchEulerDiscreteScheduler

View File

@@ -12,7 +12,7 @@ specific language governing permissions and limitations under the License.
# DPMSolverMultistepScheduler
`DPMSolverMultistepScheduler` is a multistep scheduler from [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095) by Cheng Lu, Yuhao Zhou, Fan Bao, Jianfei Chen, Chongxuan Li, and Jun Zhu.
`DPMSolverMultistep` is a multistep scheduler from [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095) by Cheng Lu, Yuhao Zhou, Fan Bao, Jianfei Chen, Chongxuan Li, and Jun Zhu.
DPMSolver (and the improved version DPMSolver++) is a fast dedicated high-order solver for diffusion ODEs with convergence order guarantee. Empirically, DPMSolver sampling with only 20 steps can generate high-quality
samples, and it can generate quite good samples even in 10 steps.

View File

@@ -37,7 +37,3 @@ Utility and helper functions for working with 🤗 Diffusers.
## make_image_grid
[[autodoc]] utils.make_image_grid
## randn_tensor
[[autodoc]] utils.torch_utils.randn_tensor

View File

@@ -1,21 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Video Processor
The [`VideoProcessor`] provides a unified API for video pipelines to prepare inputs for VAE encoding and post-processing outputs once they're decoded. The class inherits [`VaeImageProcessor`] so it includes transformations such as resizing, normalization, and conversion between PIL Image, PyTorch, and NumPy arrays.
## VideoProcessor
[[autodoc]] video_processor.VideoProcessor.preprocess_video
[[autodoc]] video_processor.VideoProcessor.postprocess_video

View File

@@ -198,81 +198,38 @@ Anything displayed on [the official Diffusers doc page](https://huggingface.co/d
Please have a look at [this page](https://github.com/huggingface/diffusers/tree/main/docs) on how to verify changes made to the documentation locally.
### 6. Contribute a community pipeline
> [!TIP]
> Read the [Community pipelines](../using-diffusers/custom_pipeline_overview#community-pipelines) guide to learn more about the difference between a GitHub and Hugging Face Hub community pipeline. If you're interested in why we have community pipelines, take a look at GitHub Issue [#841](https://github.com/huggingface/diffusers/issues/841) (basically, we can't maintain all the possible ways diffusion models can be used for inference but we also don't want to prevent the community from building them).
[Pipelines](https://huggingface.co/docs/diffusers/api/pipelines/overview) are usually the first point of contact between the Diffusers library and the user.
Pipelines are examples of how to use Diffusers [models](https://huggingface.co/docs/diffusers/api/models/overview) and [schedulers](https://huggingface.co/docs/diffusers/api/schedulers/overview).
We support two types of pipelines:
Contributing a community pipeline is a great way to share your creativity and work with the community. It lets you build on top of the [`DiffusionPipeline`] so that anyone can load and use it by setting the `custom_pipeline` parameter. This section will walk you through how to create a simple pipeline where the UNet only does a single forward pass and calls the scheduler once (a "one-step" pipeline).
- Official Pipelines
- Community Pipelines
1. Create a one_step_unet.py file for your community pipeline. This file can contain whatever package you want to use as long as it's installed by the user. Make sure you only have one pipeline class that inherits from [`DiffusionPipeline`] to load model weights and the scheduler configuration from the Hub. Add a UNet and scheduler to the `__init__` function.
Both official and community pipelines follow the same design and consist of the same type of components.
You should also add the `register_modules` function to ensure your pipeline and its components can be saved with [`~DiffusionPipeline.save_pretrained`].
Official pipelines are tested and maintained by the core maintainers of Diffusers. Their code
resides in [src/diffusers/pipelines](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines).
In contrast, community pipelines are contributed and maintained purely by the **community** and are **not** tested.
They reside in [examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community) and while they can be accessed via the [PyPI diffusers package](https://pypi.org/project/diffusers/), their code is not part of the PyPI distribution.
```py
from diffusers import DiffusionPipeline
import torch
The reason for the distinction is that the core maintainers of the Diffusers library cannot maintain and test all
possible ways diffusion models can be used for inference, but some of them may be of interest to the community.
Officially released diffusion pipelines,
such as Stable Diffusion are added to the core src/diffusers/pipelines package which ensures
high quality of maintenance, no backward-breaking code changes, and testing.
More bleeding edge pipelines should be added as community pipelines. If usage for a community pipeline is high, the pipeline can be moved to the official pipelines upon request from the community. This is one of the ways we strive to be a community-driven library.
class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
def __init__(self, unet, scheduler):
super().__init__()
To add a community pipeline, one should add a <name-of-the-community>.py file to [examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community) and adapt the [examples/community/README.md](https://github.com/huggingface/diffusers/tree/main/examples/community/README.md) to include an example of the new pipeline.
self.register_modules(unet=unet, scheduler=scheduler)
```
An example can be seen [here](https://github.com/huggingface/diffusers/pull/2400).
1. In the forward pass (which we recommend defining as `__call__`), you can add any feature you'd like. For the "one-step" pipeline, create a random image and call the UNet and scheduler once by setting `timestep=1`.
Community pipeline PRs are only checked at a superficial level and ideally they should be maintained by their original authors.
```py
from diffusers import DiffusionPipeline
import torch
class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
def __init__(self, unet, scheduler):
super().__init__()
self.register_modules(unet=unet, scheduler=scheduler)
def __call__(self):
image = torch.randn(
(1, self.unet.config.in_channels, self.unet.config.sample_size, self.unet.config.sample_size),
)
timestep = 1
model_output = self.unet(image, timestep).sample
scheduler_output = self.scheduler.step(model_output, timestep, image).prev_sample
return scheduler_output
```
Now you can run the pipeline by passing a UNet and scheduler to it or load pretrained weights if the pipeline structure is identical.
```py
from diffusers import DDPMScheduler, UNet2DModel
scheduler = DDPMScheduler()
unet = UNet2DModel()
pipeline = UnetSchedulerOneForwardPipeline(unet=unet, scheduler=scheduler)
output = pipeline()
# load pretrained weights
pipeline = UnetSchedulerOneForwardPipeline.from_pretrained("google/ddpm-cifar10-32", use_safetensors=True)
output = pipeline()
```
You can either share your pipeline as a GitHub community pipeline or Hub community pipeline.
<hfoptions id="pipeline type">
<hfoption id="GitHub pipeline">
Share your GitHub pipeline by opening a pull request on the Diffusers [repository](https://github.com/huggingface/diffusers) and add the one_step_unet.py file to the [examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community) subfolder.
</hfoption>
<hfoption id="Hub pipeline">
Share your Hub pipeline by creating a model repository on the Hub and uploading the one_step_unet.py file to it.
</hfoption>
</hfoptions>
Contributing a community pipeline is a great way to understand how Diffusers models and schedulers work. Having contributed a community pipeline is usually the first stepping stone to contributing an official pipeline to the
core package.
### 7. Contribute to training examples

View File

@@ -70,7 +70,7 @@ The following design principles are followed:
- Pipelines should be used **only** for inference.
- Pipelines should be very readable, self-explanatory, and easy to tweak.
- Pipelines should be designed to build on top of each other and be easy to integrate into higher-level APIs.
- Pipelines are **not** intended to be feature-complete user interfaces. For feature-complete user interfaces one should rather have a look at [InvokeAI](https://github.com/invoke-ai/InvokeAI), [Diffuzers](https://github.com/abhishekkrthakur/diffuzers), and [lama-cleaner](https://github.com/Sanster/lama-cleaner).
- Pipelines are **not** intended to be feature-complete user interfaces. For future complete user interfaces one should rather have a look at [InvokeAI](https://github.com/invoke-ai/InvokeAI), [Diffuzers](https://github.com/abhishekkrthakur/diffuzers), and [lama-cleaner](https://github.com/Sanster/lama-cleaner).
- Every pipeline should have one and only one way to run it via a `__call__` method. The naming of the `__call__` arguments should be shared across all pipelines.
- Pipelines should be named after the task they are intended to solve.
- In almost all cases, novel diffusion pipelines shall be implemented in a new pipeline folder/file.

View File

@@ -112,7 +112,7 @@ pip install -e ".[flax]"
These commands will link the folder you cloned the repository to and your Python library paths.
Python will now look inside the folder you cloned to in addition to the normal library paths.
For example, if your Python packages are typically installed in `~/anaconda3/envs/main/lib/python3.10/site-packages/`, Python will also search the `~/diffusers/` folder you cloned to.
For example, if your Python packages are typically installed in `~/anaconda3/envs/main/lib/python3.8/site-packages/`, Python will also search the `~/diffusers/` folder you cloned to.
<Tip warning={true}>

View File

@@ -36,7 +36,7 @@ Then load and enable the [`DeepCacheSDHelper`](https://github.com/horseee/DeepCa
image = pipe("a photo of an astronaut on a moon").images[0]
```
The `set_params` method accepts two arguments: `cache_interval` and `cache_branch_id`. `cache_interval` means the frequency of feature caching, specified as the number of steps between each cache operation. `cache_branch_id` identifies which branch of the network (ordered from the shallowest to the deepest layer) is responsible for executing the caching processes.
The `set_params` method accepts two arguments: `cache_interval` and `cache_branch_id`. `cache_interval` means the frequency of feature caching, specified as the number of steps between each cache operation. `cache_branch_id` identifies which branch of the network (ordered from the shallowest to the deepest layer) is responsible for executing the caching processes.
Opting for a lower `cache_branch_id` or a larger `cache_interval` can lead to faster inference speed at the expense of reduced image quality (ablation experiments of these two hyperparameters can be found in the [paper](https://arxiv.org/abs/2312.00858)). Once those arguments are set, use the `enable` or `disable` methods to activate or deactivate the `DeepCacheSDHelper`.
<div class="flex justify-center">

View File

@@ -12,23 +12,27 @@ specific language governing permissions and limitations under the License.
# Speed up inference
There are several ways to optimize Diffusers for inference speed, such as reducing the computational burden by lowering the data precision or using a lightweight distilled model. There are also memory-efficient attention implementations, [xFormers](xformers) and [scaled dot product attention](https://pytorch.org/docs/stable/generated/torch.nn.functional.scaled_dot_product_attention.html) in PyTorch 2.0, that reduce memory usage which also indirectly speeds up inference. Different speed optimizations can be stacked together to get the fastest inference times.
There are several ways to optimize 🤗 Diffusers for inference speed. As a general rule of thumb, we recommend using either [xFormers](xformers) or `torch.nn.functional.scaled_dot_product_attention` in PyTorch 2.0 for their memory-efficient attention.
> [!TIP]
> Optimizing for inference speed or reduced memory usage can lead to improved performance in the other category, so you should try to optimize for both whenever you can. This guide focuses on inference speed, but you can learn more about lowering memory usage in the [Reduce memory usage](memory) guide.
<Tip>
The inference times below are obtained from generating a single 512x512 image from the prompt "a photo of an astronaut riding a horse on mars" with 50 DDIM steps on a NVIDIA A100.
In many cases, optimizing for speed or memory leads to improved performance in the other, so you should try to optimize for both whenever you can. This guide focuses on inference speed, but you can learn more about preserving memory in the [Reduce memory usage](memory) guide.
| setup | latency | speed-up |
|----------|---------|----------|
| baseline | 5.27s | x1 |
| tf32 | 4.14s | x1.27 |
| fp16 | 3.51s | x1.50 |
| combined | 3.41s | x1.54 |
</Tip>
## TensorFloat-32
The results below are obtained from generating a single 512x512 image from the prompt `a photo of an astronaut riding a horse on mars` with 50 DDIM steps on a Nvidia Titan RTX, demonstrating the speed-up you can expect.
On Ampere and later CUDA devices, matrix multiplications and convolutions can use the [TensorFloat-32 (tf32)](https://blogs.nvidia.com/blog/2020/05/14/tensorfloat-32-precision-format/) mode for faster, but slightly less accurate computations. By default, PyTorch enables tf32 mode for convolutions but not matrix multiplications. Unless your network requires full float32 precision, we recommend enabling tf32 for matrix multiplications. It can significantly speed up computations with typically negligible loss in numerical accuracy.
| | latency | speed-up |
| ---------------- | ------- | ------- |
| original | 9.50s | x1 |
| fp16 | 3.61s | x2.63 |
| channels last | 3.30s | x2.88 |
| traced UNet | 3.21s | x2.96 |
| memory efficient attention | 2.63s | x3.61 |
## Use TensorFloat-32
On Ampere and later CUDA devices, matrix multiplications and convolutions can use the [TensorFloat-32 (TF32)](https://blogs.nvidia.com/blog/2020/05/14/tensorfloat-32-precision-format/) mode for faster, but slightly less accurate computations. By default, PyTorch enables TF32 mode for convolutions but not matrix multiplications. Unless your network requires full float32 precision, we recommend enabling TF32 for matrix multiplications. It can significantly speeds up computations with typically negligible loss in numerical accuracy.
```python
import torch
@@ -36,11 +40,11 @@ import torch
torch.backends.cuda.matmul.allow_tf32 = True
```
Learn more about tf32 in the [Mixed precision training](https://huggingface.co/docs/transformers/en/perf_train_gpu_one#tf32) guide.
You can learn more about TF32 in the [Mixed precision training](https://huggingface.co/docs/transformers/en/perf_train_gpu_one#tf32) guide.
## Half-precision weights
To save GPU memory and get more speed, set `torch_dtype=torch.float16` to load and run the model weights directly with half-precision weights.
To save GPU memory and get more speed, try loading and running the model weights directly in half-precision or float16:
```Python
import torch
@@ -52,76 +56,19 @@ pipe = DiffusionPipeline.from_pretrained(
use_safetensors=True,
)
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
image = pipe(prompt).images[0]
```
> [!WARNING]
> Don't use [torch.autocast](https://pytorch.org/docs/stable/amp.html#torch.autocast) in any of the pipelines as it can lead to black images and is always slower than pure float16 precision.
<Tip warning={true}>
Don't use [`torch.autocast`](https://pytorch.org/docs/stable/amp.html#torch.autocast) in any of the pipelines as it can lead to black images and is always slower than pure float16 precision.
</Tip>
## Distilled model
You could also use a distilled Stable Diffusion model and autoencoder to speed up inference. During distillation, many of the UNet's residual and attention blocks are shed to reduce the model size by 51% and improve latency on CPU/GPU by 43%. The distilled model is faster and uses less memory while generating images of comparable quality to the full Stable Diffusion model.
You could also use a distilled Stable Diffusion model and autoencoder to speed up inference. During distillation, many of the UNet's residual and attention blocks are shed to reduce the model size. The distilled model is faster and uses less memory while generating images of comparable quality to the full Stable Diffusion model.
> [!TIP]
> Read the [Open-sourcing Knowledge Distillation Code and Weights of SD-Small and SD-Tiny](https://huggingface.co/blog/sd_distillation) blog post to learn more about how knowledge distillation training works to produce a faster, smaller, and cheaper generative model.
The inference times below are obtained from generating 4 images from the prompt "a photo of an astronaut riding a horse on mars" with 25 PNDM steps on a NVIDIA A100. Each generation is repeated 3 times with the distilled Stable Diffusion v1.4 model by [Nota AI](https://hf.co/nota-ai).
| setup | latency | speed-up |
|------------------------------|---------|----------|
| baseline | 6.37s | x1 |
| distilled | 4.18s | x1.52 |
| distilled + tiny autoencoder | 3.83s | x1.66 |
Let's load the distilled Stable Diffusion model and compare it against the original Stable Diffusion model.
```py
from diffusers import StableDiffusionPipeline
import torch
distilled = StableDiffusionPipeline.from_pretrained(
"nota-ai/bk-sdm-small", torch_dtype=torch.float16, use_safetensors=True,
).to("cuda")
prompt = "a golden vase with different flowers"
generator = torch.manual_seed(2023)
image = distilled("a golden vase with different flowers", num_inference_steps=25, generator=generator).images[0]
image
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/original_sd.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">original Stable Diffusion</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/distilled_sd.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">distilled Stable Diffusion</figcaption>
</div>
</div>
### Tiny AutoEncoder
To speed inference up even more, replace the autoencoder with a [distilled version](https://huggingface.co/sayakpaul/taesdxl-diffusers) of it.
```py
import torch
from diffusers import AutoencoderTiny, StableDiffusionPipeline
distilled = StableDiffusionPipeline.from_pretrained(
"nota-ai/bk-sdm-small", torch_dtype=torch.float16, use_safetensors=True,
).to("cuda")
distilled.vae = AutoencoderTiny.from_pretrained(
"sayakpaul/taesd-diffusers", torch_dtype=torch.float16, use_safetensors=True,
).to("cuda")
prompt = "a golden vase with different flowers"
generator = torch.manual_seed(2023)
image = distilled("a golden vase with different flowers", num_inference_steps=25, generator=generator).images[0]
image
```
<div class="flex justify-center">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/distilled_sd_vae.png" />
<figcaption class="mt-2 text-center text-sm text-gray-500">distilled Stable Diffusion + Tiny AutoEncoder</figcaption>
</div>
</div>
Learn more about in the [Distilled Stable Diffusion inference](../using-diffusers/distilled_sd) guide!

View File

@@ -261,7 +261,7 @@ from dataclasses import dataclass
@dataclass
class UNet2DConditionOutput:
sample: torch.Tensor
sample: torch.FloatTensor
pipe = StableDiffusionPipeline.from_pretrained(

View File

@@ -0,0 +1,17 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Overview
Generating high-quality outputs is computationally intensive, especially during each iterative step where you go from a noisy output to a less noisy output. One of 🤗 Diffuser's goals is to make this technology widely accessible to everyone, which includes enabling fast inference on consumer and specialized hardware.
This section will cover tips and tricks - like half-precision weights and sliced attention - for optimizing inference speed and reducing memory-consumption. You'll also learn how to speed up your PyTorch code with [`torch.compile`](https://pytorch.org/tutorials/intermediate/torch_compile_tutorial.html) or [ONNX Runtime](https://onnxruntime.ai/docs/), and enable memory-efficient attention with [xFormers](https://facebookresearch.github.io/xformers/). There are also guides for running inference on specific hardware like Apple Silicon, and Intel or Habana processors.

View File

@@ -1,182 +0,0 @@
# T-GATE
[T-GATE](https://github.com/HaozheLiu-ST/T-GATE/tree/main) accelerates inference for [Stable Diffusion](../api/pipelines/stable_diffusion/overview), [PixArt](../api/pipelines/pixart), and [Latency Consistency Model](../api/pipelines/latent_consistency_models.md) pipelines by skipping the cross-attention calculation once it converges. This method doesn't require any additional training and it can speed up inference from 10-50%. T-GATE is also compatible with other optimization methods like [DeepCache](./deepcache).
Before you begin, make sure you install T-GATE.
```bash
pip install tgate
pip install -U torch diffusers transformers accelerate DeepCache
```
To use T-GATE with a pipeline, you need to use its corresponding loader.
| Pipeline | T-GATE Loader |
|---|---|
| PixArt | TgatePixArtLoader |
| Stable Diffusion XL | TgateSDXLLoader |
| Stable Diffusion XL + DeepCache | TgateSDXLDeepCacheLoader |
| Stable Diffusion | TgateSDLoader |
| Stable Diffusion + DeepCache | TgateSDDeepCacheLoader |
Next, create a `TgateLoader` with a pipeline, the gate step (the time step to stop calculating the cross attention), and the number of inference steps. Then call the `tgate` method on the pipeline with a prompt, gate step, and the number of inference steps.
Let's see how to enable this for several different pipelines.
<hfoptions id="pipelines">
<hfoption id="PixArt">
Accelerate `PixArtAlphaPipeline` with T-GATE:
```py
import torch
from diffusers import PixArtAlphaPipeline
from tgate import TgatePixArtLoader
pipe = PixArtAlphaPipeline.from_pretrained("PixArt-alpha/PixArt-XL-2-1024-MS", torch_dtype=torch.float16)
gate_step = 8
inference_step = 25
pipe = TgatePixArtLoader(
pipe,
gate_step=gate_step,
num_inference_steps=inference_step,
).to("cuda")
image = pipe.tgate(
"An alpaca made of colorful building blocks, cyberpunk.",
gate_step=gate_step,
num_inference_steps=inference_step,
).images[0]
```
</hfoption>
<hfoption id="Stable Diffusion XL">
Accelerate `StableDiffusionXLPipeline` with T-GATE:
```py
import torch
from diffusers import StableDiffusionXLPipeline
from diffusers import DPMSolverMultistepScheduler
from tgate import TgateSDXLLoader
pipe = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
variant="fp16",
use_safetensors=True,
)
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
gate_step = 10
inference_step = 25
pipe = TgateSDXLLoader(
pipe,
gate_step=gate_step,
num_inference_steps=inference_step,
).to("cuda")
image = pipe.tgate(
"Astronaut in a jungle, cold color palette, muted colors, detailed, 8k.",
gate_step=gate_step,
num_inference_steps=inference_step
).images[0]
```
</hfoption>
<hfoption id="StableDiffusionXL with DeepCache">
Accelerate `StableDiffusionXLPipeline` with [DeepCache](https://github.com/horseee/DeepCache) and T-GATE:
```py
import torch
from diffusers import StableDiffusionXLPipeline
from diffusers import DPMSolverMultistepScheduler
from tgate import TgateSDXLDeepCacheLoader
pipe = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
variant="fp16",
use_safetensors=True,
)
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
gate_step = 10
inference_step = 25
pipe = TgateSDXLDeepCacheLoader(
pipe,
cache_interval=3,
cache_branch_id=0,
).to("cuda")
image = pipe.tgate(
"Astronaut in a jungle, cold color palette, muted colors, detailed, 8k.",
gate_step=gate_step,
num_inference_steps=inference_step
).images[0]
```
</hfoption>
<hfoption id="Latent Consistency Model">
Accelerate `latent-consistency/lcm-sdxl` with T-GATE:
```py
import torch
from diffusers import StableDiffusionXLPipeline
from diffusers import UNet2DConditionModel, LCMScheduler
from diffusers import DPMSolverMultistepScheduler
from tgate import TgateSDXLLoader
unet = UNet2DConditionModel.from_pretrained(
"latent-consistency/lcm-sdxl",
torch_dtype=torch.float16,
variant="fp16",
)
pipe = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
unet=unet,
torch_dtype=torch.float16,
variant="fp16",
)
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
gate_step = 1
inference_step = 4
pipe = TgateSDXLLoader(
pipe,
gate_step=gate_step,
num_inference_steps=inference_step,
lcm=True
).to("cuda")
image = pipe.tgate(
"Astronaut in a jungle, cold color palette, muted colors, detailed, 8k.",
gate_step=gate_step,
num_inference_steps=inference_step
).images[0]
```
</hfoption>
</hfoptions>
T-GATE also supports [`StableDiffusionPipeline`] and [PixArt-alpha/PixArt-LCM-XL-2-1024-MS](https://hf.co/PixArt-alpha/PixArt-LCM-XL-2-1024-MS).
## Benchmarks
| Model | MACs | Param | Latency | Zero-shot 10K-FID on MS-COCO |
|-----------------------|----------|-----------|---------|---------------------------|
| SD-1.5 | 16.938T | 859.520M | 7.032s | 23.927 |
| SD-1.5 w/ T-GATE | 9.875T | 815.557M | 4.313s | 20.789 |
| SD-2.1 | 38.041T | 865.785M | 16.121s | 22.609 |
| SD-2.1 w/ T-GATE | 22.208T | 815.433 M | 9.878s | 19.940 |
| SD-XL | 149.438T | 2.570B | 53.187s | 24.628 |
| SD-XL w/ T-GATE | 84.438T | 2.024B | 27.932s | 22.738 |
| Pixart-Alpha | 107.031T | 611.350M | 61.502s | 38.669 |
| Pixart-Alpha w/ T-GATE | 65.318T | 462.585M | 37.867s | 35.825 |
| DeepCache (SD-XL) | 57.888T | - | 19.931s | 23.755 |
| DeepCache w/ T-GATE | 43.868T | - | 14.666s | 23.999 |
| LCM (SD-XL) | 11.955T | 2.570B | 3.805s | 25.044 |
| LCM w/ T-GATE | 11.171T | 2.024B | 3.533s | 25.028 |
| LCM (Pixart-Alpha) | 8.563T | 611.350M | 4.733s | 36.086 |
| LCM w/ T-GATE | 7.623T | 462.585M | 4.543s | 37.048 |
The latency is tested on an NVIDIA 1080TI, MACs and Params are calculated with [calflops](https://github.com/MrYxJ/calculate-flops.pytorch), and the FID is calculated with [PytorchFID](https://github.com/mseitzer/pytorch-fid).

View File

@@ -49,7 +49,7 @@ One of the simplest ways to speed up inference is to place the pipeline on a GPU
pipeline = pipeline.to("cuda")
```
To make sure you can use the same image and improve on it, use a [`Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) and set a seed for [reproducibility](./using-diffusers/reusing_seeds):
To make sure you can use the same image and improve on it, use a [`Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) and set a seed for [reproducibility](./using-diffusers/reproducibility):
```python
import torch

View File

@@ -88,7 +88,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()

View File

@@ -54,7 +54,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -84,7 +84,7 @@ Many of the basic parameters are described in the [DreamBooth](dreambooth#script
- `--freeze_model`: freezes the key and value parameters in the cross-attention layer; the default is `crossattn_kv`, but you can set it to `crossattn` to train all the parameters in the cross-attention layer
- `--concepts_list`: to learn multiple concepts, provide a path to a JSON file containing the concepts
- `--modifier_token`: a special word used to represent the learned concept
- `--initializer_token`: a special word used to initialize the embeddings of the `modifier_token`
- `--initializer_token`:
### Prior preservation loss

View File

@@ -52,76 +52,6 @@ To learn more, take a look at the [Distributed Inference with 🤗 Accelerate](h
</Tip>
### Device placement
> [!WARNING]
> This feature is experimental and its APIs might change in the future.
With Accelerate, you can use the `device_map` to determine how to distribute the models of a pipeline across multiple devices. This is useful in situations where you have more than one GPU.
For example, if you have two 8GB GPUs, then using [`~DiffusionPipeline.enable_model_cpu_offload`] may not work so well because:
* it only works on a single GPU
* a single model might not fit on a single GPU ([`~DiffusionPipeline.enable_sequential_cpu_offload`] might work but it will be extremely slow and it is also limited to a single GPU)
To make use of both GPUs, you can use the "balanced" device placement strategy which splits the models across all available GPUs.
> [!WARNING]
> Only the "balanced" strategy is supported at the moment, and we plan to support additional mapping strategies in the future.
```diff
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained(
- "runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True,
+ "runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True, device_map="balanced"
)
image = pipeline("a dog").images[0]
image
```
You can also pass a dictionary to enforce the maximum GPU memory that can be used on each device:
```diff
from diffusers import DiffusionPipeline
import torch
max_memory = {0:"1GB", 1:"1GB"}
pipeline = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
device_map="balanced",
+ max_memory=max_memory
)
image = pipeline("a dog").images[0]
image
```
If a device is not present in `max_memory`, then it will be completely ignored and will not participate in the device placement.
By default, Diffusers uses the maximum memory of all devices. If the models don't fit on the GPUs, they are offloaded to the CPU. If the CPU doesn't have enough memory, then you might see an error. In that case, you could defer to using [`~DiffusionPipeline.enable_sequential_cpu_offload`] and [`~DiffusionPipeline.enable_model_cpu_offload`].
Call [`~DiffusionPipeline.reset_device_map`] to reset the `device_map` of a pipeline. This is also necessary if you want to use methods like `to()`, [`~DiffusionPipeline.enable_sequential_cpu_offload`], and [`~DiffusionPipeline.enable_model_cpu_offload`] on a pipeline that was device-mapped.
```py
pipeline.reset_device_map()
```
Once a pipeline has been device-mapped, you can also access its device map via `hf_device_map`:
```py
print(pipeline.hf_device_map)
```
An example device map would look like so:
```bash
{'unet': 1, 'vae': 1, 'safety_checker': 0, 'text_encoder': 0}
```
## PyTorch Distributed
PyTorch supports [`DistributedDataParallel`](https://pytorch.org/docs/stable/generated/torch.nn.parallel.DistributedDataParallel.html) which enables data parallelism.

View File

@@ -67,7 +67,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -180,7 +180,7 @@ elif args.pretrained_model_name_or_path:
revision=args.revision,
use_fast=False,
)
# Load scheduler and models
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
text_encoder = text_encoder_cls.from_pretrained(
@@ -440,198 +440,6 @@ Stable Diffusion XL (SDXL) is a powerful text-to-image model that generates high
The SDXL training script is discussed in more detail in the [SDXL training](sdxl) guide.
## DeepFloyd IF
DeepFloyd IF is a cascading pixel diffusion model with three stages. The first stage generates a base image and the second and third stages progressively upscales the base image into a high-resolution 1024x1024 image. Use the [train_dreambooth_lora.py](https://github.com/huggingface/diffusers/blob/main/examples/dreambooth/train_dreambooth_lora.py) or [train_dreambooth.py](https://github.com/huggingface/diffusers/blob/main/examples/dreambooth/train_dreambooth.py) scripts to train a DeepFloyd IF model with LoRA or the full model.
DeepFloyd IF uses predicted variance, but the Diffusers training scripts uses predicted error so the trained DeepFloyd IF models are switched to a fixed variance schedule. The training scripts will update the scheduler config of the fully trained model for you. However, when you load the saved LoRA weights you must also update the pipeline's scheduler config.
```py
from diffusers import DiffusionPipeline
pipe = DiffusionPipeline.from_pretrained("DeepFloyd/IF-I-XL-v1.0", use_safetensors=True)
pipe.load_lora_weights("<lora weights path>")
# Update scheduler config to fixed variance schedule
pipe.scheduler = pipe.scheduler.__class__.from_config(pipe.scheduler.config, variance_type="fixed_small")
```
The stage 2 model requires additional validation images to upscale. You can download and use a downsized version of the training images for this.
```py
from huggingface_hub import snapshot_download
local_dir = "./dog_downsized"
snapshot_download(
"diffusers/dog-example-downsized",
local_dir=local_dir,
repo_type="dataset",
ignore_patterns=".gitattributes",
)
```
The code samples below provide a brief overview of how to train a DeepFloyd IF model with a combination of DreamBooth and LoRA. Some important parameters to note are:
* `--resolution=64`, a much smaller resolution is required because DeepFloyd IF is a pixel diffusion model and to work on uncompressed pixels, the input images must be smaller
* `--pre_compute_text_embeddings`, compute the text embeddings ahead of time to save memory because the [`~transformers.T5Model`] can take up a lot of memory
* `--tokenizer_max_length=77`, you can use a longer default text length with T5 as the text encoder but the default model encoding procedure uses a shorter text length
* `--text_encoder_use_attention_mask`, to pass the attention mask to the text encoder
<hfoptions id="IF-DreamBooth">
<hfoption id="Stage 1 LoRA DreamBooth">
Training stage 1 of DeepFloyd IF with LoRA and DreamBooth requires ~28GB of memory.
```bash
export MODEL_NAME="DeepFloyd/IF-I-XL-v1.0"
export INSTANCE_DIR="dog"
export OUTPUT_DIR="dreambooth_dog_lora"
accelerate launch train_dreambooth_lora.py \
--report_to wandb \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--output_dir=$OUTPUT_DIR \
--instance_prompt="a sks dog" \
--resolution=64 \
--train_batch_size=4 \
--gradient_accumulation_steps=1 \
--learning_rate=5e-6 \
--scale_lr \
--max_train_steps=1200 \
--validation_prompt="a sks dog" \
--validation_epochs=25 \
--checkpointing_steps=100 \
--pre_compute_text_embeddings \
--tokenizer_max_length=77 \
--text_encoder_use_attention_mask
```
</hfoption>
<hfoption id="Stage 2 LoRA DreamBooth">
For stage 2 of DeepFloyd IF with LoRA and DreamBooth, pay attention to these parameters:
* `--validation_images`, the images to upscale during validation
* `--class_labels_conditioning=timesteps`, to additionally conditional the UNet as needed in stage 2
* `--learning_rate=1e-6`, a lower learning rate is used compared to stage 1
* `--resolution=256`, the expected resolution for the upscaler
```bash
export MODEL_NAME="DeepFloyd/IF-II-L-v1.0"
export INSTANCE_DIR="dog"
export OUTPUT_DIR="dreambooth_dog_upscale"
export VALIDATION_IMAGES="dog_downsized/image_1.png dog_downsized/image_2.png dog_downsized/image_3.png dog_downsized/image_4.png"
python train_dreambooth_lora.py \
--report_to wandb \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--output_dir=$OUTPUT_DIR \
--instance_prompt="a sks dog" \
--resolution=256 \
--train_batch_size=4 \
--gradient_accumulation_steps=1 \
--learning_rate=1e-6 \
--max_train_steps=2000 \
--validation_prompt="a sks dog" \
--validation_epochs=100 \
--checkpointing_steps=500 \
--pre_compute_text_embeddings \
--tokenizer_max_length=77 \
--text_encoder_use_attention_mask \
--validation_images $VALIDATION_IMAGES \
--class_labels_conditioning=timesteps
```
</hfoption>
<hfoption id="Stage 1 DreamBooth">
For stage 1 of DeepFloyd IF with DreamBooth, pay attention to these parameters:
* `--skip_save_text_encoder`, to skip saving the full T5 text encoder with the finetuned model
* `--use_8bit_adam`, to use 8-bit Adam optimizer to save memory due to the size of the optimizer state when training the full model
* `--learning_rate=1e-7`, a really low learning rate should be used for full model training otherwise the model quality is degraded (you can use a higher learning rate with a larger batch size)
Training with 8-bit Adam and a batch size of 4, the full model can be trained with ~48GB of memory.
```bash
export MODEL_NAME="DeepFloyd/IF-I-XL-v1.0"
export INSTANCE_DIR="dog"
export OUTPUT_DIR="dreambooth_if"
accelerate launch train_dreambooth.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--output_dir=$OUTPUT_DIR \
--instance_prompt="a photo of sks dog" \
--resolution=64 \
--train_batch_size=4 \
--gradient_accumulation_steps=1 \
--learning_rate=1e-7 \
--max_train_steps=150 \
--validation_prompt "a photo of sks dog" \
--validation_steps 25 \
--text_encoder_use_attention_mask \
--tokenizer_max_length 77 \
--pre_compute_text_embeddings \
--use_8bit_adam \
--set_grads_to_none \
--skip_save_text_encoder \
--push_to_hub
```
</hfoption>
<hfoption id="Stage 2 DreamBooth">
For stage 2 of DeepFloyd IF with DreamBooth, pay attention to these parameters:
* `--learning_rate=5e-6`, use a lower learning rate with a smaller effective batch size
* `--resolution=256`, the expected resolution for the upscaler
* `--train_batch_size=2` and `--gradient_accumulation_steps=6`, to effectively train on images wiht faces requires larger batch sizes
```bash
export MODEL_NAME="DeepFloyd/IF-II-L-v1.0"
export INSTANCE_DIR="dog"
export OUTPUT_DIR="dreambooth_dog_upscale"
export VALIDATION_IMAGES="dog_downsized/image_1.png dog_downsized/image_2.png dog_downsized/image_3.png dog_downsized/image_4.png"
accelerate launch train_dreambooth.py \
--report_to wandb \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--output_dir=$OUTPUT_DIR \
--instance_prompt="a sks dog" \
--resolution=256 \
--train_batch_size=2 \
--gradient_accumulation_steps=6 \
--learning_rate=5e-6 \
--max_train_steps=2000 \
--validation_prompt="a sks dog" \
--validation_steps=150 \
--checkpointing_steps=500 \
--pre_compute_text_embeddings \
--tokenizer_max_length=77 \
--text_encoder_use_attention_mask \
--validation_images $VALIDATION_IMAGES \
--class_labels_conditioning timesteps \
--push_to_hub
```
</hfoption>
</hfoptions>
### Training tips
Training the DeepFloyd IF model can be challenging, but here are some tips that we've found helpful:
- LoRA is sufficient for training the stage 1 model because the model's low resolution makes representing finer details difficult regardless.
- For common or simple objects, you don't necessarily need to finetune the upscaler. Make sure the prompt passed to the upscaler is adjusted to remove the new token from the instance prompt. For example, if your stage 1 prompt is "a sks dog" then your stage 2 prompt should be "a dog".
- For finer details like faces, fully training the stage 2 upscaler is better than training the stage 2 model with LoRA. It also helps to use lower learning rates with larger batch sizes.
- Lower learning rates should be used to train the stage 2 model.
- The [`DDPMScheduler`] works better than the DPMSolver used in the training scripts.
## Next steps
Congratulations on training your DreamBooth model! To learn more about how to use your new model, the following guide may be helpful:

View File

@@ -51,7 +51,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -89,7 +89,7 @@ The dataset preprocessing code and training loop are found in the [`main()`](htt
As with the script parameters, a walkthrough of the training script is provided in the [Text-to-image](text2image#training-script) training guide. Instead, this guide takes a look at the InstructPix2Pix relevant parts of the script.
The script begins by modifying the [number of input channels](https://github.com/huggingface/diffusers/blob/64603389da01082055a901f2883c4810d1144edb/examples/instruct_pix2pix/train_instruct_pix2pix.py#L445) in the first convolutional layer of the UNet to account for InstructPix2Pix's additional conditioning image:
The script begins by modifing the [number of input channels](https://github.com/huggingface/diffusers/blob/64603389da01082055a901f2883c4810d1144edb/examples/instruct_pix2pix/train_instruct_pix2pix.py#L445) in the first convolutional layer of the UNet to account for InstructPix2Pix's additional conditioning image:
```py
in_channels = 8

View File

@@ -59,7 +59,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -205,7 +205,7 @@ model_pred = unet(noisy_latents, timesteps, None, added_cond_kwargs=added_cond_k
Once youve made all your changes or youre okay with the default configuration, youre ready to launch the training script! 🚀
You'll train on the [Naruto BLIP captions](https://huggingface.co/datasets/lambdalabs/naruto-blip-captions) dataset to generate your own Naruto characters, but you can also create and train on your own dataset by following the [Create a dataset for training](create_dataset) guide. Set the environment variable `DATASET_NAME` to the name of the dataset on the Hub or if you're training on your own files, set the environment variable `TRAIN_DIR` to a path to your dataset.
You'll train on the [Pokémon BLIP captions](https://huggingface.co/datasets/lambdalabs/pokemon-blip-captions) dataset to generate your own Pokémon, but you can also create and train on your own dataset by following the [Create a dataset for training](create_dataset) guide. Set the environment variable `DATASET_NAME` to the name of the dataset on the Hub or if you're training on your own files, set the environment variable `TRAIN_DIR` to a path to your dataset.
If youre training on more than one GPU, add the `--multi_gpu` parameter to the `accelerate launch` command.
@@ -219,7 +219,7 @@ To monitor training progress with Weights & Biases, add the `--report_to=wandb`
<hfoption id="prior model">
```bash
export DATASET_NAME="lambdalabs/naruto-blip-captions"
export DATASET_NAME="lambdalabs/pokemon-blip-captions"
accelerate launch --mixed_precision="fp16" train_text_to_image_prior.py \
--dataset_name=$DATASET_NAME \
@@ -232,17 +232,17 @@ accelerate launch --mixed_precision="fp16" train_text_to_image_prior.py \
--checkpoints_total_limit=3 \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--validation_prompts="A robot naruto, 4k photo" \
--validation_prompts="A robot pokemon, 4k photo" \
--report_to="wandb" \
--push_to_hub \
--output_dir="kandi2-prior-naruto-model"
--output_dir="kandi2-prior-pokemon-model"
```
</hfoption>
<hfoption id="decoder model">
```bash
export DATASET_NAME="lambdalabs/naruto-blip-captions"
export DATASET_NAME="lambdalabs/pokemon-blip-captions"
accelerate launch --mixed_precision="fp16" train_text_to_image_decoder.py \
--dataset_name=$DATASET_NAME \
@@ -256,10 +256,10 @@ accelerate launch --mixed_precision="fp16" train_text_to_image_decoder.py \
--checkpoints_total_limit=3 \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--validation_prompts="A robot naruto, 4k photo" \
--validation_prompts="A robot pokemon, 4k photo" \
--report_to="wandb" \
--push_to_hub \
--output_dir="kandi2-decoder-naruto-model"
--output_dir="kandi2-decoder-pokemon-model"
```
</hfoption>
@@ -279,7 +279,7 @@ prior_components = {"prior_" + k: v for k,v in prior_pipeline.components.items()
pipeline = AutoPipelineForText2Image.from_pretrained("kandinsky-community/kandinsky-2-2-decoder", **prior_components, torch_dtype=torch.float16)
pipe.enable_model_cpu_offload()
prompt="A robot naruto, 4k photo"
prompt="A robot pokemon, 4k photo"
image = pipeline(prompt=prompt, negative_prompt=negative_prompt).images[0]
```
@@ -299,7 +299,7 @@ import torch
pipeline = AutoPipelineForText2Image.from_pretrained("path/to/saved/model", torch_dtype=torch.float16)
pipeline.enable_model_cpu_offload()
prompt="A robot naruto, 4k photo"
prompt="A robot pokemon, 4k photo"
image = pipeline(prompt=prompt).images[0]
```
@@ -313,7 +313,7 @@ unet = UNet2DConditionModel.from_pretrained("path/to/saved/model" + "/checkpoint
pipeline = AutoPipelineForText2Image.from_pretrained("kandinsky-community/kandinsky-2-2-decoder", unet=unet, torch_dtype=torch.float16)
pipeline.enable_model_cpu_offload()
image = pipeline(prompt="A robot naruto, 4k photo").images[0]
image = pipeline(prompt="A robot pokemon, 4k photo").images[0]
```
</hfoption>

View File

@@ -53,7 +53,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -252,4 +252,4 @@ The SDXL training script is discussed in more detail in the [SDXL training](sdxl
Congratulations on distilling a LCM model! To learn more about LCM, the following may be helpful:
- Learn how to use [LCMs for inference](../using-diffusers/lcm) for text-to-image, image-to-image, and with LoRA checkpoints.
- Read the [SDXL in 4 steps with Latent Consistency LoRAs](https://huggingface.co/blog/lcm_lora) blog post to learn more about SDXL LCM-LoRA's for super fast inference, quality comparisons, benchmarks, and more.
- Read the [SDXL in 4 steps with Latent Consistency LoRAs](https://huggingface.co/blog/lcm_lora) blog post to learn more about SDXL LCM-LoRA's for super fast inference, quality comparisons, benchmarks, and more.

View File

@@ -170,7 +170,7 @@ Aside from setting up the LoRA layers, the training script is more or less the s
Once you've made all your changes or you're okay with the default configuration, you're ready to launch the training script! 🚀
Let's train on the [Naruto BLIP captions](https://huggingface.co/datasets/lambdalabs/naruto-blip-captions) dataset to generate your own Naruto characters. Set the environment variables `MODEL_NAME` and `DATASET_NAME` to the model and dataset respectively. You should also specify where to save the model in `OUTPUT_DIR`, and the name of the model to save to on the Hub with `HUB_MODEL_ID`. The script creates and saves the following files to your repository:
Let's train on the [Pokémon BLIP captions](https://huggingface.co/datasets/lambdalabs/pokemon-blip-captions) dataset to generate our own Pokémon. Set the environment variables `MODEL_NAME` and `DATASET_NAME` to the model and dataset respectively. You should also specify where to save the model in `OUTPUT_DIR`, and the name of the model to save to on the Hub with `HUB_MODEL_ID`. The script creates and saves the following files to your repository:
- saved model checkpoints
- `pytorch_lora_weights.safetensors` (the trained LoRA weights)
@@ -185,9 +185,9 @@ A full training run takes ~5 hours on a 2080 Ti GPU with 11GB of VRAM.
```bash
export MODEL_NAME="runwayml/stable-diffusion-v1-5"
export OUTPUT_DIR="/sddata/finetune/lora/naruto"
export HUB_MODEL_ID="naruto-lora"
export DATASET_NAME="lambdalabs/naruto-blip-captions"
export OUTPUT_DIR="/sddata/finetune/lora/pokemon"
export HUB_MODEL_ID="pokemon-lora"
export DATASET_NAME="lambdalabs/pokemon-blip-captions"
accelerate launch --mixed_precision="fp16" train_text_to_image_lora.py \
--pretrained_model_name_or_path=$MODEL_NAME \
@@ -208,7 +208,7 @@ accelerate launch --mixed_precision="fp16" train_text_to_image_lora.py \
--hub_model_id=${HUB_MODEL_ID} \
--report_to=wandb \
--checkpointing_steps=500 \
--validation_prompt="A naruto with blue eyes." \
--validation_prompt="A pokemon with blue eyes." \
--seed=1337
```
@@ -220,7 +220,7 @@ import torch
pipeline = AutoPipelineForText2Image.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16).to("cuda")
pipeline.load_lora_weights("path/to/lora/model", weight_name="pytorch_lora_weights.safetensors")
image = pipeline("A naruto with blue eyes").images[0]
image = pipeline("A pokemon with blue eyes").images[0]
```
## Next steps

View File

@@ -59,7 +59,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -176,7 +176,7 @@ If you want to learn more about how the training loop works, check out the [Unde
Once youve made all your changes or youre okay with the default configuration, youre ready to launch the training script! 🚀
Lets train on the [Naruto BLIP captions](https://huggingface.co/datasets/lambdalabs/naruto-blip-captions) dataset to generate your own Naruto characters. Set the environment variables `MODEL_NAME` and `DATASET_NAME` to the model and the dataset (either from the Hub or a local path). You should also specify a VAE other than the SDXL VAE (either from the Hub or a local path) with `VAE_NAME` to avoid numerical instabilities.
Lets train on the [Pokémon BLIP captions](https://huggingface.co/datasets/lambdalabs/pokemon-blip-captions) dataset to generate your own Pokémon. Set the environment variables `MODEL_NAME` and `DATASET_NAME` to the model and the dataset (either from the Hub or a local path). You should also specify a VAE other than the SDXL VAE (either from the Hub or a local path) with `VAE_NAME` to avoid numerical instabilities.
<Tip>
@@ -187,7 +187,7 @@ To monitor training progress with Weights & Biases, add the `--report_to=wandb`
```bash
export MODEL_NAME="stabilityai/stable-diffusion-xl-base-1.0"
export VAE_NAME="madebyollin/sdxl-vae-fp16-fix"
export DATASET_NAME="lambdalabs/naruto-blip-captions"
export DATASET_NAME="lambdalabs/pokemon-blip-captions"
accelerate launch train_text_to_image_sdxl.py \
--pretrained_model_name_or_path=$MODEL_NAME \
@@ -211,7 +211,7 @@ accelerate launch train_text_to_image_sdxl.py \
--validation_prompt="a cute Sundar Pichai creature" \
--validation_epochs 5 \
--checkpointing_steps=5000 \
--output_dir="sdxl-naruto-model" \
--output_dir="sdxl-pokemon-model" \
--push_to_hub
```
@@ -226,9 +226,9 @@ import torch
pipeline = DiffusionPipeline.from_pretrained("path/to/your/model", torch_dtype=torch.float16).to("cuda")
prompt = "A naruto with green eyes and red legs."
prompt = "A pokemon with green eyes and red legs."
image = pipeline(prompt, num_inference_steps=30, guidance_scale=7.5).images[0]
image.save("naruto.png")
image.save("pokemon.png")
```
</hfoption>
@@ -244,11 +244,11 @@ import torch_xla.core.xla_model as xm
device = xm.xla_device()
pipeline = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0").to(device)
prompt = "A naruto with green eyes and red legs."
prompt = "A pokemon with green eyes and red legs."
start = time()
image = pipeline(prompt, num_inference_steps=inference_steps).images[0]
print(f'Compilation time is {time()-start} sec')
image.save("naruto.png")
image.save("pokemon.png")
start = time()
image = pipeline(prompt, num_inference_steps=inference_steps).images[0]

View File

@@ -53,7 +53,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()

View File

@@ -69,7 +69,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -158,7 +158,7 @@ Once you've made all your changes or you're okay with the default configuration,
<hfoptions id="training-inference">
<hfoption id="PyTorch">
Let's train on the [Naruto BLIP captions](https://huggingface.co/datasets/lambdalabs/naruto-blip-captions) dataset to generate your own Naruto characters. Set the environment variables `MODEL_NAME` and `dataset_name` to the model and the dataset (either from the Hub or a local path). If you're training on more than one GPU, add the `--multi_gpu` parameter to the `accelerate launch` command.
Let's train on the [Pokémon BLIP captions](https://huggingface.co/datasets/lambdalabs/pokemon-blip-captions) dataset to generate your own Pokémon. Set the environment variables `MODEL_NAME` and `dataset_name` to the model and the dataset (either from the Hub or a local path). If you're training on more than one GPU, add the `--multi_gpu` parameter to the `accelerate launch` command.
<Tip>
@@ -168,7 +168,7 @@ To train on a local dataset, set the `TRAIN_DIR` and `OUTPUT_DIR` environment va
```bash
export MODEL_NAME="runwayml/stable-diffusion-v1-5"
export dataset_name="lambdalabs/naruto-blip-captions"
export dataset_name="lambdalabs/pokemon-blip-captions"
accelerate launch --mixed_precision="fp16" train_text_to_image.py \
--pretrained_model_name_or_path=$MODEL_NAME \
@@ -183,7 +183,7 @@ accelerate launch --mixed_precision="fp16" train_text_to_image.py \
--max_grad_norm=1 \
--enable_xformers_memory_efficient_attention
--lr_scheduler="constant" --lr_warmup_steps=0 \
--output_dir="sd-naruto-model" \
--output_dir="sd-pokemon-model" \
--push_to_hub
```
@@ -202,7 +202,7 @@ To train on a local dataset, set the `TRAIN_DIR` and `OUTPUT_DIR` environment va
```bash
export MODEL_NAME="runwayml/stable-diffusion-v1-5"
export dataset_name="lambdalabs/naruto-blip-captions"
export dataset_name="lambdalabs/pokemon-blip-captions"
python train_text_to_image_flax.py \
--pretrained_model_name_or_path=$MODEL_NAME \
@@ -212,7 +212,7 @@ python train_text_to_image_flax.py \
--max_train_steps=15000 \
--learning_rate=1e-05 \
--max_grad_norm=1 \
--output_dir="sd-naruto-model" \
--output_dir="sd-pokemon-model" \
--push_to_hub
```
@@ -231,7 +231,7 @@ import torch
pipeline = StableDiffusionPipeline.from_pretrained("path/to/saved_model", torch_dtype=torch.float16, use_safetensors=True).to("cuda")
image = pipeline(prompt="yoda").images[0]
image.save("yoda-naruto.png")
image.save("yoda-pokemon.png")
```
</hfoption>
@@ -246,7 +246,7 @@ from diffusers import FlaxStableDiffusionPipeline
pipeline, params = FlaxStableDiffusionPipeline.from_pretrained("path/to/saved_model", dtype=jax.numpy.bfloat16)
prompt = "yoda naruto"
prompt = "yoda pokemon"
prng_seed = jax.random.PRNGKey(0)
num_inference_steps = 50
@@ -261,7 +261,7 @@ prompt_ids = shard(prompt_ids)
images = pipeline(prompt_ids, params, prng_seed, num_inference_steps, jit=True).images
images = pipeline.numpy_to_pil(np.asarray(images.reshape((num_samples,) + images.shape[-3:])))
image.save("yoda-naruto.png")
image.save("yoda-pokemon.png")
```
</hfoption>

View File

@@ -67,7 +67,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()

View File

@@ -51,7 +51,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()

View File

@@ -53,7 +53,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -131,7 +131,7 @@ If you want to learn more about how the training loop works, check out the [Unde
Once youve made all your changes or youre okay with the default configuration, youre ready to launch the training script! 🚀
Set the `DATASET_NAME` environment variable to the dataset name from the Hub. This guide uses the [Naruto BLIP captions](https://huggingface.co/datasets/lambdalabs/naruto-blip-captions) dataset, but you can create and train on your own datasets as well (see the [Create a dataset for training](create_dataset) guide).
Set the `DATASET_NAME` environment variable to the dataset name from the Hub. This guide uses the [Pokémon BLIP captions](https://huggingface.co/datasets/lambdalabs/pokemon-blip-captions) dataset, but you can create and train on your own datasets as well (see the [Create a dataset for training](create_dataset) guide).
<Tip>
@@ -140,7 +140,7 @@ To monitor training progress with Weights & Biases, add the `--report_to=wandb`
</Tip>
```bash
export DATASET_NAME="lambdalabs/naruto-blip-captions"
export DATASET_NAME="lambdalabs/pokemon-blip-captions"
accelerate launch train_text_to_image_prior.py \
--mixed_precision="fp16" \
@@ -156,10 +156,10 @@ accelerate launch train_text_to_image_prior.py \
--checkpoints_total_limit=3 \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--validation_prompts="A robot naruto, 4k photo" \
--validation_prompts="A robot pokemon, 4k photo" \
--report_to="wandb" \
--push_to_hub \
--output_dir="wuerstchen-prior-naruto-model"
--output_dir="wuerstchen-prior-pokemon-model"
```
Once training is complete, you can use your newly trained model for inference!
@@ -171,9 +171,9 @@ from diffusers.pipelines.wuerstchen import DEFAULT_STAGE_C_TIMESTEPS
pipeline = AutoPipelineForText2Image.from_pretrained("path/to/saved/model", torch_dtype=torch.float16).to("cuda")
caption = "A cute bird naruto holding a shield"
caption = "A cute bird pokemon holding a shield"
images = pipeline(
caption,
caption,
width=1024,
height=1536,
prior_timesteps=DEFAULT_STAGE_C_TIMESTEPS,

View File

@@ -12,74 +12,75 @@ specific language governing permissions and limitations under the License.
# AutoPipeline
Diffusers provides many pipelines for basic tasks like generating images, videos, audio, and inpainting. On top of these, there are specialized pipelines for adapters and features like upscaling, super-resolution, and more. Different pipeline classes can even use the same checkpoint because they share the same pretrained model! With so many different pipelines, it can be overwhelming to know which pipeline class to use.
🤗 Diffusers is able to complete many different tasks, and you can often reuse the same pretrained weights for multiple tasks such as text-to-image, image-to-image, and inpainting. If you're new to the library and diffusion models though, it may be difficult to know which pipeline to use for a task. For example, if you're using the [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5) checkpoint for text-to-image, you might not know that you could also use it for image-to-image and inpainting by loading the checkpoint with the [`StableDiffusionImg2ImgPipeline`] and [`StableDiffusionInpaintPipeline`] classes respectively.
The [AutoPipeline](../api/pipelines/auto_pipeline) class is designed to simplify the variety of pipelines in Diffusers. It is a generic *task-first* pipeline that lets you focus on a task ([`AutoPipelineForText2Image`], [`AutoPipelineForImage2Image`], and [`AutoPipelineForInpainting`]) without needing to know the specific pipeline class. The [AutoPipeline](../api/pipelines/auto_pipeline) automatically detects the correct pipeline class to use.
The `AutoPipeline` class is designed to simplify the variety of pipelines in 🤗 Diffusers. It is a generic, *task-first* pipeline that lets you focus on the task. The `AutoPipeline` automatically detects the correct pipeline class to use, which makes it easier to load a checkpoint for a task without knowing the specific pipeline class name.
For example, let's use the [dreamlike-art/dreamlike-photoreal-2.0](https://hf.co/dreamlike-art/dreamlike-photoreal-2.0) checkpoint.
<Tip>
Under the hood, [AutoPipeline](../api/pipelines/auto_pipeline):
Take a look at the [AutoPipeline](../api/pipelines/auto_pipeline) reference to see which tasks are supported. Currently, it supports text-to-image, image-to-image, and inpainting.
1. Detects a `"stable-diffusion"` class from the [model_index.json](https://hf.co/dreamlike-art/dreamlike-photoreal-2.0/blob/main/model_index.json) file.
2. Depending on the task you're interested in, it loads the [`StableDiffusionPipeline`], [`StableDiffusionImg2ImgPipeline`], or [`StableDiffusionInpaintPipeline`]. Any parameter (`strength`, `num_inference_steps`, etc.) you would pass to these specific pipelines can also be passed to the [AutoPipeline](../api/pipelines/auto_pipeline).
</Tip>
<hfoptions id="autopipeline">
<hfoption id="text-to-image">
This tutorial shows you how to use an `AutoPipeline` to automatically infer the pipeline class to load for a specific task, given the pretrained weights.
## Choose an AutoPipeline for your task
Start by picking a checkpoint. For example, if you're interested in text-to-image with the [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5) checkpoint, use [`AutoPipelineForText2Image`]:
```py
from diffusers import AutoPipelineForText2Image
import torch
pipe_txt2img = AutoPipelineForText2Image.from_pretrained(
"dreamlike-art/dreamlike-photoreal-2.0", torch_dtype=torch.float16, use_safetensors=True
pipeline = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
prompt = "peasant and dragon combat, wood cutting style, viking era, bevel with rune"
prompt = "cinematic photo of Godzilla eating sushi with a cat in a izakaya, 35mm photograph, film, professional, 4k, highly detailed"
generator = torch.Generator(device="cpu").manual_seed(37)
image = pipe_txt2img(prompt, generator=generator).images[0]
image = pipeline(prompt, num_inference_steps=25).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-text2img.png"/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-text2img.png" alt="generated image of peasant fighting dragon in wood cutting style"/>
</div>
</hfoption>
<hfoption id="image-to-image">
Under the hood, [`AutoPipelineForText2Image`]:
1. automatically detects a `"stable-diffusion"` class from the [`model_index.json`](https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/model_index.json) file
2. loads the corresponding text-to-image [`StableDiffusionPipeline`] based on the `"stable-diffusion"` class name
Likewise, for image-to-image, [`AutoPipelineForImage2Image`] detects a `"stable-diffusion"` checkpoint from the `model_index.json` file and it'll load the corresponding [`StableDiffusionImg2ImgPipeline`] behind the scenes. You can also pass any additional arguments specific to the pipeline class such as `strength`, which determines the amount of noise or variation added to an input image:
```py
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import load_image
import torch
import requests
from PIL import Image
from io import BytesIO
pipe_img2img = AutoPipelineForImage2Image.from_pretrained(
"dreamlike-art/dreamlike-photoreal-2.0", torch_dtype=torch.float16, use_safetensors=True
pipeline = AutoPipelineForImage2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
).to("cuda")
prompt = "a portrait of a dog wearing a pearl earring"
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-text2img.png")
url = "https://upload.wikimedia.org/wikipedia/commons/thumb/0/0f/1665_Girl_with_a_Pearl_Earring.jpg/800px-1665_Girl_with_a_Pearl_Earring.jpg"
prompt = "cinematic photo of Godzilla eating burgers with a cat in a fast food restaurant, 35mm photograph, film, professional, 4k, highly detailed"
generator = torch.Generator(device="cpu").manual_seed(53)
image = pipe_img2img(prompt, image=init_image, generator=generator).images[0]
response = requests.get(url)
image = Image.open(BytesIO(response.content)).convert("RGB")
image.thumbnail((768, 768))
image = pipeline(prompt, image, num_inference_steps=200, strength=0.75, guidance_scale=10.5).images[0]
image
```
Notice how the [dreamlike-art/dreamlike-photoreal-2.0](https://hf.co/dreamlike-art/dreamlike-photoreal-2.0) checkpoint is used for both text-to-image and image-to-image tasks? To save memory and avoid loading the checkpoint twice, use the [`~DiffusionPipeline.from_pipe`] method.
```py
pipe_img2img = AutoPipelineForImage2Image.from_pipe(pipe_txt2img).to("cuda")
image = pipeline(prompt, image=init_image, generator=generator).images[0]
image
```
You can learn more about the [`~DiffusionPipeline.from_pipe`] method in the [Reuse a pipeline](../using-diffusers/loading#reuse-a-pipeline) guide.
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-img2img.png"/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-img2img.png" alt="generated image of a vermeer portrait of a dog wearing a pearl earring"/>
</div>
</hfoption>
<hfoption id="inpainting">
And if you want to do inpainting, then [`AutoPipelineForInpainting`] loads the underlying [`StableDiffusionInpaintPipeline`] class in the same way:
```py
from diffusers import AutoPipelineForInpainting
@@ -90,27 +91,22 @@ pipeline = AutoPipelineForInpainting.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-img2img.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-mask.png")
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
prompt = "cinematic photo of a owl, 35mm photograph, film, professional, 4k, highly detailed"
generator = torch.Generator(device="cpu").manual_seed(38)
image = pipeline(prompt, image=init_image, mask_image=mask_image, generator=generator, strength=0.4).images[0]
init_image = load_image(img_url).convert("RGB")
mask_image = load_image(mask_url).convert("RGB")
prompt = "A majestic tiger sitting on a bench"
image = pipeline(prompt, image=init_image, mask_image=mask_image, num_inference_steps=50, strength=0.80).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-inpaint.png"/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-inpaint.png" alt="generated image of a tiger sitting on a bench"/>
</div>
</hfoption>
</hfoptions>
## Unsupported checkpoints
The [AutoPipeline](../api/pipelines/auto_pipeline) supports [Stable Diffusion](../api/pipelines/stable_diffusion/overview), [Stable Diffusion XL](../api/pipelines/stable_diffusion/stable_diffusion_xl), [ControlNet](../api/pipelines/controlnet), [Kandinsky 2.1](../api/pipelines/kandinsky.md), [Kandinsky 2.2](../api/pipelines/kandinsky_v22), and [DeepFloyd IF](../api/pipelines/deepfloyd_if) checkpoints.
If you try to load an unsupported checkpoint, you'll get an error.
If you try to load an unsupported checkpoint, it'll throw an error:
```py
from diffusers import AutoPipelineForImage2Image
@@ -121,3 +117,54 @@ pipeline = AutoPipelineForImage2Image.from_pretrained(
)
"ValueError: AutoPipeline can't find a pipeline linked to ShapEImg2ImgPipeline for None"
```
## Use multiple pipelines
For some workflows or if you're loading many pipelines, it is more memory-efficient to reuse the same components from a checkpoint instead of reloading them which would unnecessarily consume additional memory. For example, if you're using a checkpoint for text-to-image and you want to use it again for image-to-image, use the [`~AutoPipelineForImage2Image.from_pipe`] method. This method creates a new pipeline from the components of a previously loaded pipeline at no additional memory cost.
The [`~AutoPipelineForImage2Image.from_pipe`] method detects the original pipeline class and maps it to the new pipeline class corresponding to the task you want to do. For example, if you load a `"stable-diffusion"` class pipeline for text-to-image:
```py
from diffusers import AutoPipelineForText2Image, AutoPipelineForImage2Image
import torch
pipeline_text2img = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
)
print(type(pipeline_text2img))
"<class 'diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline'>"
```
Then [`~AutoPipelineForImage2Image.from_pipe`] maps the original `"stable-diffusion"` pipeline class to [`StableDiffusionImg2ImgPipeline`]:
```py
pipeline_img2img = AutoPipelineForImage2Image.from_pipe(pipeline_text2img)
print(type(pipeline_img2img))
"<class 'diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_img2img.StableDiffusionImg2ImgPipeline'>"
```
If you passed an optional argument - like disabling the safety checker - to the original pipeline, this argument is also passed on to the new pipeline:
```py
from diffusers import AutoPipelineForText2Image, AutoPipelineForImage2Image
import torch
pipeline_text2img = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
requires_safety_checker=False,
).to("cuda")
pipeline_img2img = AutoPipelineForImage2Image.from_pipe(pipeline_text2img)
print(pipeline_img2img.config.requires_safety_checker)
"False"
```
You can overwrite any of the arguments and even configuration from the original pipeline if you want to change the behavior of the new pipeline. For example, to turn the safety checker back on and add the `strength` argument:
```py
pipeline_img2img = AutoPipelineForImage2Image.from_pipe(pipeline_text2img, requires_safety_checker=True, strength=0.3)
print(pipeline_img2img.config.requires_safety_checker)
"True"
```

View File

@@ -260,7 +260,7 @@ Then, you'll need a way to evaluate the model. For evaluation, you can use the [
... # The default pipeline output type is `List[PIL.Image]`
... images = pipeline(
... batch_size=config.eval_batch_size,
... generator=torch.Generator(device='cpu').manual_seed(config.seed), # Use a separate torch generator to avoid rewinding the random state of the main training loop
... generator=torch.manual_seed(config.seed),
... ).images
... # Make a grid out of the images

View File

@@ -14,17 +14,19 @@ specific language governing permissions and limitations under the License.
# Load LoRAs for inference
There are many adapter types (with [LoRAs](https://huggingface.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora) being the most popular) trained in different styles to achieve different effects. You can even combine multiple adapters to create new and unique images.
There are many adapters (with LoRAs being the most common type) trained in different styles to achieve different effects. You can even combine multiple adapters to create new and unique images. With the 🤗 [PEFT](https://huggingface.co/docs/peft/index) integration in 🤗 Diffusers, it is really easy to load and manage adapters for inference. In this guide, you'll learn how to use different adapters with [Stable Diffusion XL (SDXL)](../api/pipelines/stable_diffusion/stable_diffusion_xl) for inference.
In this tutorial, you'll learn how to easily load and manage adapters for inference with the 🤗 [PEFT](https://huggingface.co/docs/peft/index) integration in 🤗 Diffusers. You'll use LoRA as the main adapter technique, so you'll see the terms LoRA and adapter used interchangeably.
Throughout this guide, you'll use LoRA as the main adapter technique, so we'll use the terms LoRA and adapter interchangeably. You should have some familiarity with LoRA, and if you don't, we welcome you to check out the [LoRA guide](https://huggingface.co/docs/peft/conceptual_guides/lora).
Let's first install all the required libraries.
```bash
!pip install -q transformers accelerate peft diffusers
!pip install -q transformers accelerate
!pip install peft
!pip install diffusers
```
Now, load a pipeline with a [Stable Diffusion XL (SDXL)](../api/pipelines/stable_diffusion/stable_diffusion_xl) checkpoint:
Now, let's load a pipeline with a SDXL checkpoint:
```python
from diffusers import DiffusionPipeline
@@ -34,18 +36,21 @@ pipe_id = "stabilityai/stable-diffusion-xl-base-1.0"
pipe = DiffusionPipeline.from_pretrained(pipe_id, torch_dtype=torch.float16).to("cuda")
```
Next, load a [CiroN2022/toy-face](https://huggingface.co/CiroN2022/toy-face) adapter with the [`~diffusers.loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights`] method. With the 🤗 PEFT integration, you can assign a specific `adapter_name` to the checkpoint, which let's you easily switch between different LoRA checkpoints. Let's call this adapter `"toy"`.
Next, load a LoRA checkpoint with the [`~diffusers.loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights`] method.
With the 🤗 PEFT integration, you can assign a specific `adapter_name` to the checkpoint, which let's you easily switch between different LoRA checkpoints. Let's call this adapter `"toy"`.
```python
pipe.load_lora_weights("CiroN2022/toy-face", weight_name="toy_face_sdxl.safetensors", adapter_name="toy")
```
Make sure to include the token `toy_face` in the prompt and then you can perform inference:
And then perform inference:
```python
prompt = "toy_face of a hacker with a hoodie"
lora_scale = 0.9
lora_scale= 0.9
image = pipe(
prompt, num_inference_steps=30, cross_attention_kwargs={"scale": lora_scale}, generator=torch.manual_seed(0)
).images[0]
@@ -54,16 +59,17 @@ image
![toy-face](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_8_1.png)
With the `adapter_name` parameter, it is really easy to use another adapter for inference! Load the [nerijs/pixel-art-xl](https://huggingface.co/nerijs/pixel-art-xl) adapter that has been fine-tuned to generate pixel art images and call it `"pixel"`.
The pipeline automatically sets the first loaded adapter (`"toy"`) as the active adapter, but you can activate the `"pixel"` adapter with the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method:
With the `adapter_name` parameter, it is really easy to use another adapter for inference! Load the [nerijs/pixel-art-xl](https://huggingface.co/nerijs/pixel-art-xl) adapter that has been fine-tuned to generate pixel art images, and let's call it `"pixel"`.
The pipeline automatically sets the first loaded adapter (`"toy"`) as the active adapter. But you can activate the `"pixel"` adapter with the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method as shown below:
```python
pipe.load_lora_weights("nerijs/pixel-art-xl", weight_name="pixel-art-xl.safetensors", adapter_name="pixel")
pipe.set_adapters("pixel")
```
Make sure you include the token `pixel art` in your prompt to generate a pixel art image:
Let's now generate an image with the second adapter and check the result:
```python
prompt = "a hacker with a hoodie, pixel art"
@@ -75,25 +81,29 @@ image
![pixel-art](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_12_1.png)
## Merge adapters
## Combine multiple adapters
You can also merge different adapter checkpoints for inference to blend their styles together.
You can also perform multi-adapter inference where you combine different adapter checkpoints for inference.
Once again, use the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method to activate the `pixel` and `toy` adapters and specify the weights for how they should be merged.
Once again, use the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method to activate two LoRA checkpoints and specify the weight for how the checkpoints should be combined.
```python
pipe.set_adapters(["pixel", "toy"], adapter_weights=[0.5, 1.0])
```
Now that we have set these two adapters, let's generate an image from the combined adapters!
<Tip>
LoRA checkpoints in the diffusion community are almost always obtained with [DreamBooth](https://huggingface.co/docs/diffusers/main/en/training/dreambooth). DreamBooth training often relies on "trigger" words in the input text prompts in order for the generation results to look as expected. When you combine multiple LoRA checkpoints, it's important to ensure the trigger words for the corresponding LoRA checkpoints are present in the input text prompts.
</Tip>
Remember to use the trigger words for [CiroN2022/toy-face](https://hf.co/CiroN2022/toy-face) and [nerijs/pixel-art-xl](https://hf.co/nerijs/pixel-art-xl) (these are found in their repositories) in the prompt to generate an image.
The trigger words for [CiroN2022/toy-face](https://hf.co/CiroN2022/toy-face) and [nerijs/pixel-art-xl](https://hf.co/nerijs/pixel-art-xl) are found in their repositories.
```python
# Notice how the prompt is constructed.
prompt = "toy_face of a hacker with a hoodie, pixel art"
image = pipe(
prompt, num_inference_steps=30, cross_attention_kwargs={"scale": 1.0}, generator=torch.manual_seed(0)
@@ -103,95 +113,43 @@ image
![toy-face-pixel-art](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_16_1.png)
Impressive! As you can see, the model generated an image that mixed the characteristics of both adapters.
Impressive! As you can see, the model was able to generate an image that mixes the characteristics of both adapters.
> [!TIP]
> Through its PEFT integration, Diffusers also offers more efficient merging methods which you can learn about in the [Merge LoRAs](../using-diffusers/merge_loras) guide!
To return to only using one adapter, use the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method to activate the `"toy"` adapter:
If you want to go back to using only one adapter, use the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method to activate the `"toy"` adapter:
```python
# First, set the adapter.
pipe.set_adapters("toy")
# Then, run inference.
prompt = "toy_face of a hacker with a hoodie"
lora_scale = 0.9
lora_scale= 0.9
image = pipe(
prompt, num_inference_steps=30, cross_attention_kwargs={"scale": lora_scale}, generator=torch.manual_seed(0)
).images[0]
image
```
Or to disable all adapters entirely, use the [`~diffusers.loaders.UNet2DConditionLoadersMixin.disable_lora`] method to return the base model.
![toy-face-again](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_18_1.png)
If you want to switch to only the base model, disable all LoRAs with the [`~diffusers.loaders.UNet2DConditionLoadersMixin.disable_lora`] method.
```python
pipe.disable_lora()
prompt = "toy_face of a hacker with a hoodie"
lora_scale= 0.9
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
image
```
![no-lora](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_20_1.png)
### Customize adapters strength
For even more customization, you can control how strongly the adapter affects each part of the pipeline. For this, pass a dictionary with the control strengths (called "scales") to [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`].
## Monitoring active adapters
For example, here's how you can turn on the adapter for the `down` parts, but turn it off for the `mid` and `up` parts:
```python
pipe.enable_lora() # enable lora again, after we disabled it above
prompt = "toy_face of a hacker with a hoodie, pixel art"
adapter_weight_scales = { "unet": { "down": 1, "mid": 0, "up": 0} }
pipe.set_adapters("pixel", adapter_weight_scales)
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
image
```
![block-lora-text-and-down](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_block_down.png)
Let's see how turning off the `down` part and turning on the `mid` and `up` part respectively changes the image.
```python
adapter_weight_scales = { "unet": { "down": 0, "mid": 1, "up": 0} }
pipe.set_adapters("pixel", adapter_weight_scales)
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
image
```
![block-lora-text-and-mid](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_block_mid.png)
```python
adapter_weight_scales = { "unet": { "down": 0, "mid": 0, "up": 1} }
pipe.set_adapters("pixel", adapter_weight_scales)
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
image
```
![block-lora-text-and-up](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_block_up.png)
Looks cool!
This is a really powerful feature. You can use it to control the adapter strengths down to per-transformer level. And you can even use it for multiple adapters.
```python
adapter_weight_scales_toy = 0.5
adapter_weight_scales_pixel = {
"unet": {
"down": 0.9, # all transformers in the down-part will use scale 0.9
# "mid" # because, in this example, "mid" is not given, all transformers in the mid part will use the default scale 1.0
"up": {
"block_0": 0.6, # all 3 transformers in the 0th block in the up-part will use scale 0.6
"block_1": [0.4, 0.8, 1.0], # the 3 transformers in the 1st block in the up-part will use scales 0.4, 0.8 and 1.0 respectively
}
}
}
pipe.set_adapters(["toy", "pixel"], [adapter_weight_scales_toy, adapter_weight_scales_pixel])
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
image
```
![block-lora-mixed](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_block_mixed.png)
## Manage active adapters
You have attached multiple adapters in this tutorial, and if you're feeling a bit lost on what adapters have been attached to the pipeline's components, use the [`~diffusers.loaders.LoraLoaderMixin.get_active_adapters`] method to check the list of active adapters:
You have attached multiple adapters in this tutorial, and if you're feeling a bit lost on what adapters have been attached to the pipeline's components, you can easily check the list of active adapters using the [`~diffusers.loaders.LoraLoaderMixin.get_active_adapters`] method:
```py
active_adapters = pipe.get_active_adapters()
@@ -206,3 +164,78 @@ list_adapters_component_wise = pipe.get_list_adapters()
list_adapters_component_wise
{"text_encoder": ["toy", "pixel"], "unet": ["toy", "pixel"], "text_encoder_2": ["toy", "pixel"]}
```
## Compatibility with `torch.compile`
If you want to compile your model with `torch.compile` make sure to first fuse the LoRA weights into the base model and unload them.
```diff
pipe.load_lora_weights("nerijs/pixel-art-xl", weight_name="pixel-art-xl.safetensors", adapter_name="pixel")
pipe.load_lora_weights("CiroN2022/toy-face", weight_name="toy_face_sdxl.safetensors", adapter_name="toy")
pipe.set_adapters(["pixel", "toy"], adapter_weights=[0.5, 1.0])
# Fuses the LoRAs into the Unet
pipe.fuse_lora()
pipe.unload_lora_weights()
+ pipe.unet.to(memory_format=torch.channels_last)
+ pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
prompt = "toy_face of a hacker with a hoodie, pixel art"
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
```
> [!TIP]
> You can refer to the `torch.compile()` section [here](https://huggingface.co/docs/diffusers/main/en/optimization/torch2.0#torchcompile) and [here](https://huggingface.co/docs/diffusers/main/en/tutorials/fast_diffusion#torchcompile) for more elaborate examples.
## Fusing adapters into the model
You can use PEFT to easily fuse/unfuse multiple adapters directly into the model weights (both UNet and text encoder) using the [`~diffusers.loaders.LoraLoaderMixin.fuse_lora`] method, which can lead to a speed-up in inference and lower VRAM usage.
```py
pipe.load_lora_weights("nerijs/pixel-art-xl", weight_name="pixel-art-xl.safetensors", adapter_name="pixel")
pipe.load_lora_weights("CiroN2022/toy-face", weight_name="toy_face_sdxl.safetensors", adapter_name="toy")
pipe.set_adapters(["pixel", "toy"], adapter_weights=[0.5, 1.0])
# Fuses the LoRAs into the Unet
pipe.fuse_lora()
prompt = "toy_face of a hacker with a hoodie, pixel art"
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
# Gets the Unet back to the original state
pipe.unfuse_lora()
```
You can also fuse some adapters using `adapter_names` for faster generation:
```py
pipe.load_lora_weights("nerijs/pixel-art-xl", weight_name="pixel-art-xl.safetensors", adapter_name="pixel")
pipe.load_lora_weights("CiroN2022/toy-face", weight_name="toy_face_sdxl.safetensors", adapter_name="toy")
pipe.set_adapters(["pixel"], adapter_weights=[0.5, 1.0])
# Fuses the LoRAs into the Unet
pipe.fuse_lora(adapter_names=["pixel"])
prompt = "a hacker with a hoodie, pixel art"
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
# Gets the Unet back to the original state
pipe.unfuse_lora()
# Fuse all adapters
pipe.fuse_lora(adapter_names=["pixel", "toy"])
prompt = "toy_face of a hacker with a hoodie, pixel art"
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
```
## Saving a pipeline after fusing the adapters
To properly save a pipeline after it's been loaded with the adapters, it should be serialized like so:
```python
pipe.fuse_lora(lora_scale=1.0)
pipe.unload_lora_weights()
pipe.save_pretrained("path-to-pipeline")
```

View File

@@ -12,93 +12,27 @@ specific language governing permissions and limitations under the License.
# Pipeline callbacks
The denoising loop of a pipeline can be modified with custom defined functions using the `callback_on_step_end` parameter. The callback function is executed at the end of each step, and modifies the pipeline attributes and variables for the next step. This is really useful for *dynamically* adjusting certain pipeline attributes or modifying tensor variables. This versatility allows for interesting use-cases such as changing the prompt embeddings at each timestep, assigning different weights to the prompt embeddings, and editing the guidance scale. With callbacks, you can implement new features without modifying the underlying code!
The denoising loop of a pipeline can be modified with custom defined functions using the `callback_on_step_end` parameter. This can be really useful for *dynamically* adjusting certain pipeline attributes, or modifying tensor variables. The flexibility of callbacks opens up some interesting use-cases such as changing the prompt embeddings at each timestep, assigning different weights to the prompt embeddings, and editing the guidance scale.
> [!TIP]
> 🤗 Diffusers currently only supports `callback_on_step_end`, but feel free to open a [feature request](https://github.com/huggingface/diffusers/issues/new/choose) if you have a cool use-case and require a callback function with a different execution point!
This guide will show you how to use the `callback_on_step_end` parameter to disable classifier-free guidance (CFG) after 40% of the inference steps to save compute with minimal cost to performance.
This guide will demonstrate how callbacks work by a few features you can implement with them.
The callback function should have the following arguments:
## Official callbacks
We provide a list of callbacks you can plug into an existing pipeline and modify the denoising loop. This is the current list of official callbacks:
- `SDCFGCutoffCallback`: Disables the CFG after a certain number of steps for all SD 1.5 pipelines, including text-to-image, image-to-image, inpaint, and controlnet.
- `SDXLCFGCutoffCallback`: Disables the CFG after a certain number of steps for all SDXL pipelines, including text-to-image, image-to-image, inpaint, and controlnet.
- `IPAdapterScaleCutoffCallback`: Disables the IP Adapter after a certain number of steps for all pipelines supporting IP-Adapter.
> [!TIP]
> If you want to add a new official callback, feel free to open a [feature request](https://github.com/huggingface/diffusers/issues/new/choose) or [submit a PR](https://huggingface.co/docs/diffusers/main/en/conceptual/contribution#how-to-open-a-pr).
To set up a callback, you need to specify the number of denoising steps after which the callback comes into effect. You can do so by using either one of these two arguments
- `cutoff_step_ratio`: Float number with the ratio of the steps.
- `cutoff_step_index`: Integer number with the exact number of the step.
```python
import torch
from diffusers import DPMSolverMultistepScheduler, StableDiffusionXLPipeline
from diffusers.callbacks import SDXLCFGCutoffCallback
callback = SDXLCFGCutoffCallback(cutoff_step_ratio=0.4)
# can also be used with cutoff_step_index
# callback = SDXLCFGCutoffCallback(cutoff_step_ratio=None, cutoff_step_index=10)
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
variant="fp16",
).to("cuda")
pipeline.scheduler = DPMSolverMultistepScheduler.from_config(pipeline.scheduler.config, use_karras_sigmas=True)
prompt = "a sports car at the road, best quality, high quality, high detail, 8k resolution"
generator = torch.Generator(device="cpu").manual_seed(2628670641)
out = pipeline(
prompt=prompt,
negative_prompt="",
guidance_scale=6.5,
num_inference_steps=25,
generator=generator,
callback_on_step_end=callback,
)
out.images[0].save("official_callback.png")
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/without_cfg_callback.png" alt="generated image of a sports car at the road" />
<figcaption class="mt-2 text-center text-sm text-gray-500">without SDXLCFGCutoffCallback</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/with_cfg_callback.png" alt="generated image of a a sports car at the road with cfg callback" />
<figcaption class="mt-2 text-center text-sm text-gray-500">with SDXLCFGCutoffCallback</figcaption>
</div>
</div>
## Dynamic classifier-free guidance
Dynamic classifier-free guidance (CFG) is a feature that allows you to disable CFG after a certain number of inference steps which can help you save compute with minimal cost to performance. The callback function for this should have the following arguments:
- `pipeline` (or the pipeline instance) provides access to important properties such as `num_timesteps` and `guidance_scale`. You can modify these properties by updating the underlying attributes. For this example, you'll disable CFG by setting `pipeline._guidance_scale=0.0`.
- `step_index` and `timestep` tell you where you are in the denoising loop. Use `step_index` to turn off CFG after reaching 40% of `num_timesteps`.
- `callback_kwargs` is a dict that contains tensor variables you can modify during the denoising loop. It only includes variables specified in the `callback_on_step_end_tensor_inputs` argument, which is passed to the pipeline's `__call__` method. Different pipelines may use different sets of variables, so please check a pipeline's `_callback_tensor_inputs` attribute for the list of variables you can modify. Some common variables include `latents` and `prompt_embeds`. For this function, change the batch size of `prompt_embeds` after setting `guidance_scale=0.0` in order for it to work properly.
* `pipe` (or the pipeline instance) provides access to useful properties such as `num_timesteps` and `guidance_scale`. You can modify these properties by updating the underlying attributes. For this example, you'll disable CFG by setting `pipe._guidance_scale=0.0`.
* `step_index` and `timestep` tell you where you are in the denoising loop. Use `step_index` to turn off CFG after reaching 40% of `num_timesteps`.
* `callback_kwargs` is a dict that contains tensor variables you can modify during the denoising loop. It only includes variables specified in the `callback_on_step_end_tensor_inputs` argument, which is passed to the pipeline's `__call__` method. Different pipelines may use different sets of variables, so please check a pipeline's `_callback_tensor_inputs` attribute for the list of variables you can modify. Some common variables include `latents` and `prompt_embeds`. For this function, change the batch size of `prompt_embeds` after setting `guidance_scale=0.0` in order for it to work properly.
Your callback function should look something like this:
```python
def callback_dynamic_cfg(pipe, step_index, timestep, callback_kwargs):
# adjust the batch_size of prompt_embeds according to guidance_scale
if step_index == int(pipeline.num_timesteps * 0.4):
if step_index == int(pipe.num_timesteps * 0.4):
prompt_embeds = callback_kwargs["prompt_embeds"]
prompt_embeds = prompt_embeds.chunk(2)[-1]
# update guidance_scale and prompt_embeds
pipeline._guidance_scale = 0.0
pipe._guidance_scale = 0.0
callback_kwargs["prompt_embeds"] = prompt_embeds
return callback_kwargs
```
@@ -109,134 +43,58 @@ Now, you can pass the callback function to the `callback_on_step_end` parameter
import torch
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipeline = pipeline.to("cuda")
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
generator = torch.Generator(device="cuda").manual_seed(1)
out = pipeline(
prompt,
generator=generator,
callback_on_step_end=callback_dynamic_cfg,
callback_on_step_end_tensor_inputs=['prompt_embeds']
)
out = pipe(prompt, generator=generator, callback_on_step_end=callback_dynamic_cfg, callback_on_step_end_tensor_inputs=['prompt_embeds'])
out.images[0].save("out_custom_cfg.png")
```
The callback function is executed at the end of each denoising step, and modifies the pipeline attributes and tensor variables for the next denoising step.
With callbacks, you can implement features such as dynamic CFG without having to modify the underlying code at all!
<Tip>
🤗 Diffusers currently only supports `callback_on_step_end`, but feel free to open a [feature request](https://github.com/huggingface/diffusers/issues/new/choose) if you have a cool use-case and require a callback function with a different execution point!
</Tip>
## Interrupt the diffusion process
> [!TIP]
> The interruption callback is supported for text-to-image, image-to-image, and inpainting for the [StableDiffusionPipeline](../api/pipelines/stable_diffusion/overview) and [StableDiffusionXLPipeline](../api/pipelines/stable_diffusion/stable_diffusion_xl).
Interrupting the diffusion process is particularly useful when building UIs that work with Diffusers because it allows users to stop the generation process if they're unhappy with the intermediate results. You can incorporate this into your pipeline with a callback.
Stopping the diffusion process early is useful when building UIs that work with Diffusers because it allows users to stop the generation process if they're unhappy with the intermediate results. You can incorporate this into your pipeline with a callback.
<Tip>
This callback function should take the following arguments: `pipeline`, `i`, `t`, and `callback_kwargs` (this must be returned). Set the pipeline's `_interrupt` attribute to `True` to stop the diffusion process after a certain number of steps. You are also free to implement your own custom stopping logic inside the callback.
The interruption callback is supported for text-to-image, image-to-image, and inpainting for the [StableDiffusionPipeline](../api/pipelines/stable_diffusion/overview) and [StableDiffusionXLPipeline](../api/pipelines/stable_diffusion/stable_diffusion_xl).
</Tip>
This callback function should take the following arguments: `pipe`, `i`, `t`, and `callback_kwargs` (this must be returned). Set the pipeline's `_interrupt` attribute to `True` to stop the diffusion process after a certain number of steps. You are also free to implement your own custom stopping logic inside the callback.
In this example, the diffusion process is stopped after 10 steps even though `num_inference_steps` is set to 50.
```python
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
pipeline.enable_model_cpu_offload()
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
pipe.enable_model_cpu_offload()
num_inference_steps = 50
def interrupt_callback(pipeline, i, t, callback_kwargs):
def interrupt_callback(pipe, i, t, callback_kwargs):
stop_idx = 10
if i == stop_idx:
pipeline._interrupt = True
pipe._interrupt = True
return callback_kwargs
pipeline(
pipe(
"A photo of a cat",
num_inference_steps=num_inference_steps,
callback_on_step_end=interrupt_callback,
)
```
## Display image after each generation step
> [!TIP]
> This tip was contributed by [asomoza](https://github.com/asomoza).
Display an image after each generation step by accessing and converting the latents after each step into an image. The latent space is compressed to 128x128, so the images are also 128x128 which is useful for a quick preview.
1. Use the function below to convert the SDXL latents (4 channels) to RGB tensors (3 channels) as explained in the [Explaining the SDXL latent space](https://huggingface.co/blog/TimothyAlexisVass/explaining-the-sdxl-latent-space) blog post.
```py
def latents_to_rgb(latents):
weights = (
(60, -60, 25, -70),
(60, -5, 15, -50),
(60, 10, -5, -35)
)
weights_tensor = torch.t(torch.tensor(weights, dtype=latents.dtype).to(latents.device))
biases_tensor = torch.tensor((150, 140, 130), dtype=latents.dtype).to(latents.device)
rgb_tensor = torch.einsum("...lxy,lr -> ...rxy", latents, weights_tensor) + biases_tensor.unsqueeze(-1).unsqueeze(-1)
image_array = rgb_tensor.clamp(0, 255)[0].byte().cpu().numpy()
image_array = image_array.transpose(1, 2, 0)
return Image.fromarray(image_array)
```
2. Create a function to decode and save the latents into an image.
```py
def decode_tensors(pipe, step, timestep, callback_kwargs):
latents = callback_kwargs["latents"]
image = latents_to_rgb(latents)
image.save(f"{step}.png")
return callback_kwargs
```
3. Pass the `decode_tensors` function to the `callback_on_step_end` parameter to decode the tensors after each step. You also need to specify what you want to modify in the `callback_on_step_end_tensor_inputs` parameter, which in this case are the latents.
```py
from diffusers import AutoPipelineForText2Image
import torch
from PIL import Image
pipeline = AutoPipelineForText2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
variant="fp16",
use_safetensors=True
).to("cuda")
image = pipeline(
prompt="A croissant shaped like a cute bear.",
negative_prompt="Deformed, ugly, bad anatomy",
callback_on_step_end=decode_tensors,
callback_on_step_end_tensor_inputs=["latents"],
).images[0]
```
<div class="flex gap-4 justify-center">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/tips_step_0.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">step 0</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/tips_step_19.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">step 19
</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/tips_step_29.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">step 29</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/tips_step_39.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">step 39</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/tips_step_49.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">step 49</figcaption>
</div>
</div>

View File

@@ -0,0 +1,184 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Contribute a community pipeline
<Tip>
💡 Take a look at GitHub Issue [#841](https://github.com/huggingface/diffusers/issues/841) for more context about why we're adding community pipelines to help everyone easily share their work without being slowed down.
</Tip>
Community pipelines allow you to add any additional features you'd like on top of the [`DiffusionPipeline`]. The main benefit of building on top of the `DiffusionPipeline` is anyone can load and use your pipeline by only adding one more argument, making it super easy for the community to access.
This guide will show you how to create a community pipeline and explain how they work. To keep things simple, you'll create a "one-step" pipeline where the `UNet` does a single forward pass and calls the scheduler once.
## Initialize the pipeline
You should start by creating a `one_step_unet.py` file for your community pipeline. In this file, create a pipeline class that inherits from the [`DiffusionPipeline`] to be able to load model weights and the scheduler configuration from the Hub. The one-step pipeline needs a `UNet` and a scheduler, so you'll need to add these as arguments to the `__init__` function:
```python
from diffusers import DiffusionPipeline
import torch
class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
def __init__(self, unet, scheduler):
super().__init__()
```
To ensure your pipeline and its components (`unet` and `scheduler`) can be saved with [`~DiffusionPipeline.save_pretrained`], add them to the `register_modules` function:
```diff
from diffusers import DiffusionPipeline
import torch
class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
def __init__(self, unet, scheduler):
super().__init__()
+ self.register_modules(unet=unet, scheduler=scheduler)
```
Cool, the `__init__` step is done and you can move to the forward pass now! 🔥
## Define the forward pass
In the forward pass, which we recommend defining as `__call__`, you have complete creative freedom to add whatever feature you'd like. For our amazing one-step pipeline, create a random image and only call the `unet` and `scheduler` once by setting `timestep=1`:
```diff
from diffusers import DiffusionPipeline
import torch
class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
def __init__(self, unet, scheduler):
super().__init__()
self.register_modules(unet=unet, scheduler=scheduler)
+ def __call__(self):
+ image = torch.randn(
+ (1, self.unet.config.in_channels, self.unet.config.sample_size, self.unet.config.sample_size),
+ )
+ timestep = 1
+ model_output = self.unet(image, timestep).sample
+ scheduler_output = self.scheduler.step(model_output, timestep, image).prev_sample
+ return scheduler_output
```
That's it! 🚀 You can now run this pipeline by passing a `unet` and `scheduler` to it:
```python
from diffusers import DDPMScheduler, UNet2DModel
scheduler = DDPMScheduler()
unet = UNet2DModel()
pipeline = UnetSchedulerOneForwardPipeline(unet=unet, scheduler=scheduler)
output = pipeline()
```
But what's even better is you can load pre-existing weights into the pipeline if the pipeline structure is identical. For example, you can load the [`google/ddpm-cifar10-32`](https://huggingface.co/google/ddpm-cifar10-32) weights into the one-step pipeline:
```python
pipeline = UnetSchedulerOneForwardPipeline.from_pretrained("google/ddpm-cifar10-32", use_safetensors=True)
output = pipeline()
```
## Share your pipeline
Open a Pull Request on the 🧨 Diffusers [repository](https://github.com/huggingface/diffusers) to add your awesome pipeline in `one_step_unet.py` to the [examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community) subfolder.
Once it is merged, anyone with `diffusers >= 0.4.0` installed can use this pipeline magically 🪄 by specifying it in the `custom_pipeline` argument:
```python
from diffusers import DiffusionPipeline
pipe = DiffusionPipeline.from_pretrained(
"google/ddpm-cifar10-32", custom_pipeline="one_step_unet", use_safetensors=True
)
pipe()
```
Another way to share your community pipeline is to upload the `one_step_unet.py` file directly to your preferred [model repository](https://huggingface.co/docs/hub/models-uploading) on the Hub. Instead of specifying the `one_step_unet.py` file, pass the model repository id to the `custom_pipeline` argument:
```python
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"google/ddpm-cifar10-32", custom_pipeline="stevhliu/one_step_unet", use_safetensors=True
)
```
Take a look at the following table to compare the two sharing workflows to help you decide the best option for you:
| | GitHub community pipeline | HF Hub community pipeline |
|----------------|------------------------------------------------------------------------------------------------------------------|-------------------------------------------------------------------------------------------|
| usage | same | same |
| review process | open a Pull Request on GitHub and undergo a review process from the Diffusers team before merging; may be slower | upload directly to a Hub repository without any review; this is the fastest workflow |
| visibility | included in the official Diffusers repository and documentation | included on your HF Hub profile and relies on your own usage/promotion to gain visibility |
<Tip>
💡 You can use whatever package you want in your community pipeline file - as long as the user has it installed, everything will work fine. Make sure you have one and only one pipeline class that inherits from `DiffusionPipeline` because this is automatically detected.
</Tip>
## How do community pipelines work?
A community pipeline is a class that inherits from [`DiffusionPipeline`] which means:
- It can be loaded with the [`custom_pipeline`] argument.
- The model weights and scheduler configuration are loaded from [`pretrained_model_name_or_path`].
- The code that implements a feature in the community pipeline is defined in a `pipeline.py` file.
Sometimes you can't load all the pipeline components weights from an official repository. In this case, the other components should be passed directly to the pipeline:
```python
from diffusers import DiffusionPipeline
from transformers import CLIPImageProcessor, CLIPModel
model_id = "CompVis/stable-diffusion-v1-4"
clip_model_id = "laion/CLIP-ViT-B-32-laion2B-s34B-b79K"
feature_extractor = CLIPImageProcessor.from_pretrained(clip_model_id)
clip_model = CLIPModel.from_pretrained(clip_model_id, torch_dtype=torch.float16)
pipeline = DiffusionPipeline.from_pretrained(
model_id,
custom_pipeline="clip_guided_stable_diffusion",
clip_model=clip_model,
feature_extractor=feature_extractor,
scheduler=scheduler,
torch_dtype=torch.float16,
use_safetensors=True,
)
```
The magic behind community pipelines is contained in the following code. It allows the community pipeline to be loaded from GitHub or the Hub, and it'll be available to all 🧨 Diffusers packages.
```python
# 2. Load the pipeline class, if using custom module then load it from the Hub
# if we load from explicit class, let's use it
if custom_pipeline is not None:
pipeline_class = get_class_from_dynamic_module(
custom_pipeline, module_file=CUSTOM_PIPELINE_FILE_NAME, cache_dir=custom_pipeline
)
elif cls != DiffusionPipeline:
pipeline_class = cls
else:
diffusers_module = importlib.import_module(cls.__module__.split(".")[0])
pipeline_class = getattr(diffusers_module, config_dict["_class_name"])
```

View File

@@ -0,0 +1,58 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Control image brightness
The Stable Diffusion pipeline is mediocre at generating images that are either very bright or dark as explained in the [Common Diffusion Noise Schedules and Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) paper. The solutions proposed in the paper are currently implemented in the [`DDIMScheduler`] which you can use to improve the lighting in your images.
<Tip>
💡 Take a look at the paper linked above for more details about the proposed solutions!
</Tip>
One of the solutions is to train a model with *v prediction* and *v loss*. Add the following flag to the [`train_text_to_image.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image.py) or [`train_text_to_image_lora.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image_lora.py) scripts to enable `v_prediction`:
```bash
--prediction_type="v_prediction"
```
For example, let's use the [`ptx0/pseudo-journey-v2`](https://huggingface.co/ptx0/pseudo-journey-v2) checkpoint which has been finetuned with `v_prediction`.
Next, configure the following parameters in the [`DDIMScheduler`]:
1. `rescale_betas_zero_snr=True`, rescales the noise schedule to zero terminal signal-to-noise ratio (SNR)
2. `timestep_spacing="trailing"`, starts sampling from the last timestep
```py
from diffusers import DiffusionPipeline, DDIMScheduler
pipeline = DiffusionPipeline.from_pretrained("ptx0/pseudo-journey-v2", use_safetensors=True)
# switch the scheduler in the pipeline to use the DDIMScheduler
pipeline.scheduler = DDIMScheduler.from_config(
pipeline.scheduler.config, rescale_betas_zero_snr=True, timestep_spacing="trailing"
)
pipeline.to("cuda")
```
Finally, in your call to the pipeline, set `guidance_rescale` to prevent overexposure:
```py
prompt = "A lion in galaxies, spirals, nebulae, stars, smoke, iridescent, intricate detail, octane render, 8k"
image = pipeline(prompt, guidance_rescale=0.7).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/zero_snr.png"/>
</div>

View File

@@ -429,27 +429,6 @@ image = pipe(
make_image_grid([original_image, canny_image, image], rows=1, cols=3)
```
<Tip>
You can use a refiner model with `StableDiffusionXLControlNetPipeline` to improve image quality, just like you can with a regular `StableDiffusionXLPipeline`.
See the [Refine image quality](./sdxl#refine-image-quality) section to learn how to use the refiner model.
Make sure to use `StableDiffusionXLControlNetPipeline` and pass `image` and `controlnet_conditioning_scale`.
```py
base = StableDiffusionXLControlNetPipeline(...)
image = base(
prompt=prompt,
controlnet_conditioning_scale=0.5,
image=canny_image,
num_inference_steps=40,
denoising_end=0.8,
output_type="latent",
).images
# rest exactly as with StableDiffusionXLPipeline
```
</Tip>
## MultiControlNet
<Tip>

Some files were not shown because too many files have changed in this diff Show More