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Author SHA1 Message Date
sayakpaul
1f67e4e7f6 add a test for checking effective custom gc. 2025-01-29 11:04:12 +05:30
731 changed files with 11015 additions and 75281 deletions

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@@ -1,38 +0,0 @@
name: "\U0001F31F Remote VAE"
description: Feedback for remote VAE pilot
labels: [ "Remote VAE" ]
body:
- type: textarea
id: positive
validations:
required: true
attributes:
label: Did you like the remote VAE solution?
description: |
If you liked it, we would appreciate it if you could elaborate what you liked.
- type: textarea
id: feedback
validations:
required: true
attributes:
label: What can be improved about the current solution?
description: |
Let us know the things you would like to see improved. Note that we will work optimizing the solution once the pilot is over and we have usage.
- type: textarea
id: others
validations:
required: true
attributes:
label: What other VAEs you would like to see if the pilot goes well?
description: |
Provide a list of the VAEs you would like to see in the future if the pilot goes well.
- type: textarea
id: additional-info
attributes:
label: Notify the members of the team
description: |
Tag the following folks when submitting this feedback: @hlky @sayakpaul

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@@ -38,7 +38,6 @@ jobs:
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install pandas peft
python -m uv pip uninstall transformers && python -m uv pip install transformers==4.48.0
- name: Environment
run: |
python utils/print_env.py

View File

@@ -414,16 +414,10 @@ jobs:
config:
- backend: "bitsandbytes"
test_location: "bnb"
additional_deps: ["peft"]
- backend: "gguf"
test_location: "gguf"
additional_deps: ["peft"]
- backend: "torchao"
test_location: "torchao"
additional_deps: []
- backend: "optimum_quanto"
test_location: "quanto"
additional_deps: []
runs-on:
group: aws-g6e-xlarge-plus
container:
@@ -441,9 +435,6 @@ jobs:
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install -U ${{ matrix.config.backend }}
if [ "${{ join(matrix.config.additional_deps, ' ') }}" != "" ]; then
python -m uv pip install ${{ join(matrix.config.additional_deps, ' ') }}
fi
python -m uv pip install pytest-reportlog
- name: Environment
run: |

View File

@@ -1,17 +0,0 @@
name: PR Style Bot
on:
issue_comment:
types: [created]
permissions:
contents: write
pull-requests: write
jobs:
style:
uses: huggingface/huggingface_hub/.github/workflows/style-bot-action.yml@main
with:
python_quality_dependencies: "[quality]"
secrets:
bot_token: ${{ secrets.GITHUB_TOKEN }}

View File

@@ -2,7 +2,8 @@ name: Fast tests for PRs
on:
pull_request:
branches: [main]
branches:
- main
paths:
- "src/diffusers/**.py"
- "benchmarks/**.py"
@@ -63,7 +64,6 @@ jobs:
run: |
python utils/check_copies.py
python utils/check_dummies.py
python utils/check_support_list.py
make deps_table_check_updated
- name: Check if failure
if: ${{ failure() }}
@@ -120,8 +120,7 @@ jobs:
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git --no-deps
python -m uv pip install accelerate
- name: Environment
run: |

View File

@@ -1,296 +0,0 @@
name: Fast GPU Tests on PR
on:
pull_request:
branches: main
paths:
- "src/diffusers/models/modeling_utils.py"
- "src/diffusers/models/model_loading_utils.py"
- "src/diffusers/pipelines/pipeline_utils.py"
- "src/diffusers/pipeline_loading_utils.py"
- "src/diffusers/loaders/lora_base.py"
- "src/diffusers/loaders/lora_pipeline.py"
- "src/diffusers/loaders/peft.py"
- "tests/pipelines/test_pipelines_common.py"
- "tests/models/test_modeling_common.py"
workflow_dispatch:
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
env:
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
HF_HUB_ENABLE_HF_TRANSFER: 1
PYTEST_TIMEOUT: 600
PIPELINE_USAGE_CUTOFF: 1000000000 # set high cutoff so that only always-test pipelines run
jobs:
check_code_quality:
runs-on: ubuntu-22.04
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check quality
run: make quality
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Quality check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make style && make quality'" >> $GITHUB_STEP_SUMMARY
check_repository_consistency:
needs: check_code_quality
runs-on: ubuntu-22.04
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check repo consistency
run: |
python utils/check_copies.py
python utils/check_dummies.py
python utils/check_support_list.py
make deps_table_check_updated
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Repo consistency check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make fix-copies'" >> $GITHUB_STEP_SUMMARY
setup_torch_cuda_pipeline_matrix:
needs: [check_code_quality, check_repository_consistency]
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
- name: Environment
run: |
python utils/print_env.py
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
run: |
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
echo $matrix
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: test-pipelines.json
path: reports
torch_pipelines_cuda_tests:
name: Torch Pipelines CUDA Tests
needs: setup_torch_cuda_pipeline_matrix
strategy:
fail-fast: false
max-parallel: 8
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
- name: Environment
run: |
python utils/print_env.py
- name: Extract tests
id: extract_tests
run: |
pattern=$(python utils/extract_tests_from_mixin.py --type pipeline)
echo "$pattern" > /tmp/test_pattern.txt
echo "pattern_file=/tmp/test_pattern.txt" >> $GITHUB_OUTPUT
- name: PyTorch CUDA checkpoint tests on Ubuntu
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
if [ "${{ matrix.module }}" = "ip_adapters" ]; then
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
tests/pipelines/${{ matrix.module }}
else
pattern=$(cat ${{ steps.extract_tests.outputs.pattern_file }})
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx and $pattern" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
tests/pipelines/${{ matrix.module }}
fi
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: pipeline_${{ matrix.module }}_test_reports
path: reports
torch_cuda_tests:
name: Torch CUDA Tests
needs: [check_code_quality, check_repository_consistency]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
defaults:
run:
shell: bash
strategy:
fail-fast: false
max-parallel: 2
matrix:
module: [models, schedulers, lora, others]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
- name: Environment
run: |
python utils/print_env.py
- name: Extract tests
id: extract_tests
run: |
pattern=$(python utils/extract_tests_from_mixin.py --type ${{ matrix.module }})
echo "$pattern" > /tmp/test_pattern.txt
echo "pattern_file=/tmp/test_pattern.txt" >> $GITHUB_OUTPUT
- name: Run PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
pattern=$(cat ${{ steps.extract_tests.outputs.pattern_file }})
if [ -z "$pattern" ]; then
python -m pytest -n 1 -sv --max-worker-restart=0 --dist=loadfile -k "not Flax and not Onnx" tests/${{ matrix.module }} \
--make-reports=tests_torch_cuda_${{ matrix.module }}
else
python -m pytest -n 1 -sv --max-worker-restart=0 --dist=loadfile -k "not Flax and not Onnx and $pattern" tests/${{ matrix.module }} \
--make-reports=tests_torch_cuda_${{ matrix.module }}
fi
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_cuda_${{ matrix.module }}_stats.txt
cat reports/tests_torch_cuda_${{ matrix.module }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: torch_cuda_test_reports_${{ matrix.module }}
path: reports
run_examples_tests:
name: Examples PyTorch CUDA tests on Ubuntu
needs: [check_code_quality, check_repository_consistency]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
python -m uv pip install -e [quality,test,training]
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install timm
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=examples_torch_cuda examples/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/examples_torch_cuda_stats.txt
cat reports/examples_torch_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: examples_test_reports
path: reports

View File

@@ -349,6 +349,7 @@ jobs:
container:
image: diffusers/diffusers-pytorch-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -358,6 +359,7 @@ jobs:
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"

View File

@@ -7,8 +7,8 @@ on:
default: 'diffusers/diffusers-pytorch-cuda'
description: 'Name of the Docker image'
required: true
pr_number:
description: 'PR number to test on'
branch:
description: 'PR Branch to test on'
required: true
test:
description: 'Tests to run (e.g.: `tests/models`).'
@@ -43,8 +43,8 @@ jobs:
exit 1
fi
if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines|lora) ]]; then
echo "Error: The input string must contain either 'models', 'pipelines', or 'lora' after 'tests/'."
if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines) ]]; then
echo "Error: The input string must contain either 'models' or 'pipelines' after 'tests/'."
exit 1
fi
@@ -53,13 +53,13 @@ jobs:
exit 1
fi
echo "$PY_TEST"
shell: bash -e {0}
- name: Checkout PR branch
uses: actions/checkout@v4
with:
ref: refs/pull/${{ inputs.pr_number }}/head
ref: ${{ github.event.inputs.branch }}
repository: ${{ github.event.pull_request.head.repo.full_name }}
- name: Install pytest
run: |

View File

@@ -13,6 +13,3 @@ jobs:
fetch-depth: 0
- name: Secret Scanning
uses: trufflesecurity/trufflehog@main
with:
extra_args: --results=verified,unknown

View File

@@ -28,9 +28,9 @@ ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m uv pip install --no-cache-dir \
torch \
torchvision \
torchaudio\
torch==2.1.2 \
torchvision==0.16.2 \
torchaudio==2.1.2 \
onnxruntime \
--extra-index-url https://download.pytorch.org/whl/cpu && \
python3 -m uv pip install --no-cache-dir \

View File

@@ -76,16 +76,6 @@
- local: advanced_inference/outpaint
title: Outpainting
title: Advanced inference
- sections:
- local: hybrid_inference/overview
title: Overview
- local: hybrid_inference/vae_decode
title: VAE Decode
- local: hybrid_inference/vae_encode
title: VAE Encode
- local: hybrid_inference/api_reference
title: API Reference
title: Hybrid Inference
- sections:
- local: using-diffusers/cogvideox
title: CogVideoX
@@ -99,8 +89,6 @@
title: Kandinsky
- local: using-diffusers/ip_adapter
title: IP-Adapter
- local: using-diffusers/omnigen
title: OmniGen
- local: using-diffusers/pag
title: PAG
- local: using-diffusers/controlnet
@@ -175,8 +163,6 @@
title: gguf
- local: quantization/torchao
title: torchao
- local: quantization/quanto
title: quanto
title: Quantization Methods
- sections:
- local: optimization/fp16
@@ -265,23 +251,19 @@
sections:
- local: api/models/overview
title: Overview
- local: api/models/auto_model
title: AutoModel
- sections:
- local: api/models/controlnet
title: ControlNetModel
- local: api/models/controlnet_union
title: ControlNetUnionModel
- local: api/models/controlnet_flux
title: FluxControlNetModel
- local: api/models/controlnet_hunyuandit
title: HunyuanDiT2DControlNetModel
- local: api/models/controlnet_sana
title: SanaControlNetModel
- local: api/models/controlnet_sd3
title: SD3ControlNetModel
- local: api/models/controlnet_sparsectrl
title: SparseControlNetModel
- local: api/models/controlnet_union
title: ControlNetUnionModel
title: ControlNets
- sections:
- local: api/models/allegro_transformer3d
@@ -294,16 +276,10 @@
title: ConsisIDTransformer3DModel
- local: api/models/cogview3plus_transformer2d
title: CogView3PlusTransformer2DModel
- local: api/models/cogview4_transformer2d
title: CogView4Transformer2DModel
- local: api/models/dit_transformer2d
title: DiTTransformer2DModel
- local: api/models/easyanimate_transformer3d
title: EasyAnimateTransformer3DModel
- local: api/models/flux_transformer
title: FluxTransformer2DModel
- local: api/models/hidream_image_transformer
title: HiDreamImageTransformer2DModel
- local: api/models/hunyuan_transformer2d
title: HunyuanDiT2DModel
- local: api/models/hunyuan_video_transformer_3d
@@ -312,14 +288,10 @@
title: LatteTransformer3DModel
- local: api/models/lumina_nextdit2d
title: LuminaNextDiT2DModel
- local: api/models/lumina2_transformer2d
title: Lumina2Transformer2DModel
- local: api/models/ltx_video_transformer3d
title: LTXVideoTransformer3DModel
- local: api/models/mochi_transformer3d
title: MochiTransformer3DModel
- local: api/models/omnigen_transformer
title: OmniGenTransformer2DModel
- local: api/models/pixart_transformer2d
title: PixArtTransformer2DModel
- local: api/models/prior_transformer
@@ -334,8 +306,6 @@
title: Transformer2DModel
- local: api/models/transformer_temporal
title: TransformerTemporalModel
- local: api/models/wan_transformer_3d
title: WanTransformer3DModel
title: Transformers
- sections:
- local: api/models/stable_cascade_unet
@@ -364,12 +334,8 @@
title: AutoencoderKLHunyuanVideo
- local: api/models/autoencoderkl_ltx_video
title: AutoencoderKLLTXVideo
- local: api/models/autoencoderkl_magvit
title: AutoencoderKLMagvit
- local: api/models/autoencoderkl_mochi
title: AutoencoderKLMochi
- local: api/models/autoencoder_kl_wan
title: AutoencoderKLWan
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_dc
@@ -410,8 +376,6 @@
title: CogVideoX
- local: api/pipelines/cogview3
title: CogView3
- local: api/pipelines/cogview4
title: CogView4
- local: api/pipelines/consisid
title: ConsisID
- local: api/pipelines/consistency_models
@@ -426,8 +390,6 @@
title: ControlNet with Stable Diffusion 3
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/controlnet_sana
title: ControlNet-Sana
- local: api/pipelines/controlnetxs
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
@@ -446,14 +408,10 @@
title: DiffEdit
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/easyanimate
title: EasyAnimate
- local: api/pipelines/flux
title: Flux
- local: api/pipelines/control_flux_inpaint
title: FluxControlInpaint
- local: api/pipelines/hidream
title: HiDream-I1
- local: api/pipelines/hunyuandit
title: Hunyuan-DiT
- local: api/pipelines/hunyuan_video
@@ -480,8 +438,6 @@
title: LEDITS++
- local: api/pipelines/ltx_video
title: LTXVideo
- local: api/pipelines/lumina2
title: Lumina 2.0
- local: api/pipelines/lumina
title: Lumina-T2X
- local: api/pipelines/marigold
@@ -492,8 +448,6 @@
title: MultiDiffusion
- local: api/pipelines/musicldm
title: MusicLDM
- local: api/pipelines/omnigen
title: OmniGen
- local: api/pipelines/pag
title: PAG
- local: api/pipelines/paint_by_example
@@ -506,8 +460,6 @@
title: PixArt-Σ
- local: api/pipelines/sana
title: Sana
- local: api/pipelines/sana_sprint
title: Sana Sprint
- local: api/pipelines/self_attention_guidance
title: Self-Attention Guidance
- local: api/pipelines/semantic_stable_diffusion
@@ -568,8 +520,6 @@
title: UniDiffuser
- local: api/pipelines/value_guided_sampling
title: Value-guided sampling
- local: api/pipelines/wan
title: Wan
- local: api/pipelines/wuerstchen
title: Wuerstchen
title: Pipelines
@@ -579,10 +529,6 @@
title: Overview
- local: api/schedulers/cm_stochastic_iterative
title: CMStochasticIterativeScheduler
- local: api/schedulers/ddim_cogvideox
title: CogVideoXDDIMScheduler
- local: api/schedulers/multistep_dpm_solver_cogvideox
title: CogVideoXDPMScheduler
- local: api/schedulers/consistency_decoder
title: ConsistencyDecoderScheduler
- local: api/schedulers/cosine_dpm

View File

@@ -25,16 +25,3 @@ Customized activation functions for supporting various models in 🤗 Diffusers.
## ApproximateGELU
[[autodoc]] models.activations.ApproximateGELU
## SwiGLU
[[autodoc]] models.activations.SwiGLU
## FP32SiLU
[[autodoc]] models.activations.FP32SiLU
## LinearActivation
[[autodoc]] models.activations.LinearActivation

View File

@@ -147,20 +147,3 @@ An attention processor is a class for applying different types of attention mech
## XLAFlashAttnProcessor2_0
[[autodoc]] models.attention_processor.XLAFlashAttnProcessor2_0
## XFormersJointAttnProcessor
[[autodoc]] models.attention_processor.XFormersJointAttnProcessor
## IPAdapterXFormersAttnProcessor
[[autodoc]] models.attention_processor.IPAdapterXFormersAttnProcessor
## FluxIPAdapterJointAttnProcessor2_0
[[autodoc]] models.attention_processor.FluxIPAdapterJointAttnProcessor2_0
## XLAFluxFlashAttnProcessor2_0
[[autodoc]] models.attention_processor.XLAFluxFlashAttnProcessor2_0

View File

@@ -38,33 +38,6 @@ config = PyramidAttentionBroadcastConfig(
pipe.transformer.enable_cache(config)
```
## Faster Cache
[FasterCache](https://huggingface.co/papers/2410.19355) from Zhengyao Lv, Chenyang Si, Junhao Song, Zhenyu Yang, Yu Qiao, Ziwei Liu, Kwan-Yee K. Wong.
FasterCache is a method that speeds up inference in diffusion transformers by:
- Reusing attention states between successive inference steps, due to high similarity between them
- Skipping unconditional branch prediction used in classifier-free guidance by revealing redundancies between unconditional and conditional branch outputs for the same timestep, and therefore approximating the unconditional branch output using the conditional branch output
```python
import torch
from diffusers import CogVideoXPipeline, FasterCacheConfig
pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
pipe.to("cuda")
config = FasterCacheConfig(
spatial_attention_block_skip_range=2,
spatial_attention_timestep_skip_range=(-1, 681),
current_timestep_callback=lambda: pipe.current_timestep,
attention_weight_callback=lambda _: 0.3,
unconditional_batch_skip_range=5,
unconditional_batch_timestep_skip_range=(-1, 781),
tensor_format="BFCHW",
)
pipe.transformer.enable_cache(config)
```
### CacheMixin
[[autodoc]] CacheMixin
@@ -74,9 +47,3 @@ pipe.transformer.enable_cache(config)
[[autodoc]] PyramidAttentionBroadcastConfig
[[autodoc]] apply_pyramid_attention_broadcast
### FasterCacheConfig
[[autodoc]] FasterCacheConfig
[[autodoc]] apply_faster_cache

View File

@@ -20,11 +20,6 @@ LoRA is a fast and lightweight training method that inserts and trains a signifi
- [`FluxLoraLoaderMixin`] provides similar functions for [Flux](https://huggingface.co/docs/diffusers/main/en/api/pipelines/flux).
- [`CogVideoXLoraLoaderMixin`] provides similar functions for [CogVideoX](https://huggingface.co/docs/diffusers/main/en/api/pipelines/cogvideox).
- [`Mochi1LoraLoaderMixin`] provides similar functions for [Mochi](https://huggingface.co/docs/diffusers/main/en/api/pipelines/mochi).
- [`AuraFlowLoraLoaderMixin`] provides similar functions for [AuraFlow](https://huggingface.co/fal/AuraFlow).
- [`LTXVideoLoraLoaderMixin`] provides similar functions for [LTX-Video](https://huggingface.co/docs/diffusers/main/en/api/pipelines/ltx_video).
- [`SanaLoraLoaderMixin`] provides similar functions for [Sana](https://huggingface.co/docs/diffusers/main/en/api/pipelines/sana).
- [`HunyuanVideoLoraLoaderMixin`] provides similar functions for [HunyuanVideo](https://huggingface.co/docs/diffusers/main/en/api/pipelines/hunyuan_video).
- [`Lumina2LoraLoaderMixin`] provides similar functions for [Lumina2](https://huggingface.co/docs/diffusers/main/en/api/pipelines/lumina2).
- [`AmusedLoraLoaderMixin`] is for the [`AmusedPipeline`].
- [`LoraBaseMixin`] provides a base class with several utility methods to fuse, unfuse, unload, LoRAs and more.
@@ -57,25 +52,6 @@ To learn more about how to load LoRA weights, see the [LoRA](../../using-diffuse
## Mochi1LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.Mochi1LoraLoaderMixin
## AuraFlowLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.AuraFlowLoraLoaderMixin
## LTXVideoLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.LTXVideoLoraLoaderMixin
## SanaLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.SanaLoraLoaderMixin
## HunyuanVideoLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.HunyuanVideoLoraLoaderMixin
## Lumina2LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.Lumina2LoraLoaderMixin
## AmusedLoraLoaderMixin

View File

@@ -1,29 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# AutoModel
The `AutoModel` is designed to make it easy to load a checkpoint without needing to know the specific model class. `AutoModel` automatically retrieves the correct model class from the checkpoint `config.json` file.
```python
from diffusers import AutoModel, AutoPipelineForText2Image
unet = AutoModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", subfolder="unet")
pipe = AutoPipelineForText2Image.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", unet=unet)
```
## AutoModel
[[autodoc]] AutoModel
- all
- from_pretrained

View File

@@ -1,32 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLWan
The 3D variational autoencoder (VAE) model with KL loss used in [Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLWan
vae = AutoencoderKLWan.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", subfolder="vae", torch_dtype=torch.float32)
```
## AutoencoderKLWan
[[autodoc]] AutoencoderKLWan
- decode
- all
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

View File

@@ -18,7 +18,7 @@ The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLAllegro
vae = AutoencoderKLAllegro.from_pretrained("rhymes-ai/Allegro", subfolder="vae", torch_dtype=torch.float32).to("cuda")
vae = AutoencoderKLCogVideoX.from_pretrained("rhymes-ai/Allegro", subfolder="vae", torch_dtype=torch.float32).to("cuda")
```
## AutoencoderKLAllegro

View File

@@ -1,37 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLMagvit
The 3D variational autoencoder (VAE) model with KL loss used in [EasyAnimate](https://github.com/aigc-apps/EasyAnimate) was introduced by Alibaba PAI.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLMagvit
vae = AutoencoderKLMagvit.from_pretrained("alibaba-pai/EasyAnimateV5.1-12b-zh", subfolder="vae", torch_dtype=torch.float16).to("cuda")
```
## AutoencoderKLMagvit
[[autodoc]] AutoencoderKLMagvit
- decode
- encode
- all
## AutoencoderKLOutput
[[autodoc]] models.autoencoders.autoencoder_kl.AutoencoderKLOutput
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

View File

@@ -1,30 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# CogView4Transformer2DModel
A Diffusion Transformer model for 2D data from [CogView4]()
The model can be loaded with the following code snippet.
```python
from diffusers import CogView4Transformer2DModel
transformer = CogView4Transformer2DModel.from_pretrained("THUDM/CogView4-6B", subfolder="transformer", torch_dtype=torch.bfloat16).to("cuda")
```
## CogView4Transformer2DModel
[[autodoc]] CogView4Transformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,29 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# SanaControlNetModel
The ControlNet model was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, Maneesh Agrawala. It provides a greater degree of control over text-to-image generation by conditioning the model on additional inputs such as edge maps, depth maps, segmentation maps, and keypoints for pose detection.
The abstract from the paper is:
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
This model was contributed by [ishan24](https://huggingface.co/ishan24). ❤️
The original codebase can be found at [NVlabs/Sana](https://github.com/NVlabs/Sana), and you can find official ControlNet checkpoints on [Efficient-Large-Model's](https://huggingface.co/Efficient-Large-Model) Hub profile.
## SanaControlNetModel
[[autodoc]] SanaControlNetModel
## SanaControlNetOutput
[[autodoc]] models.controlnets.controlnet_sana.SanaControlNetOutput

View File

@@ -1,30 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# EasyAnimateTransformer3DModel
A Diffusion Transformer model for 3D data from [EasyAnimate](https://github.com/aigc-apps/EasyAnimate) was introduced by Alibaba PAI.
The model can be loaded with the following code snippet.
```python
from diffusers import EasyAnimateTransformer3DModel
transformer = EasyAnimateTransformer3DModel.from_pretrained("alibaba-pai/EasyAnimateV5.1-12b-zh", subfolder="transformer", torch_dtype=torch.float16).to("cuda")
```
## EasyAnimateTransformer3DModel
[[autodoc]] EasyAnimateTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# HiDreamImageTransformer2DModel
A Transformer model for image-like data from [HiDream-I1](https://huggingface.co/HiDream-ai).
The model can be loaded with the following code snippet.
```python
from diffusers import HiDreamImageTransformer2DModel
transformer = HiDreamImageTransformer2DModel.from_pretrained("HiDream-ai/HiDream-I1-Full", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## HiDreamImageTransformer2DModel
[[autodoc]] HiDreamImageTransformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# Lumina2Transformer2DModel
A Diffusion Transformer model for 3D video-like data was introduced in [Lumina Image 2.0](https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0) by Alpha-VLLM.
The model can be loaded with the following code snippet.
```python
from diffusers import Lumina2Transformer2DModel
transformer = Lumina2Transformer2DModel.from_pretrained("Alpha-VLLM/Lumina-Image-2.0", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## Lumina2Transformer2DModel
[[autodoc]] Lumina2Transformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,30 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# OmniGenTransformer2DModel
A Transformer model that accepts multimodal instructions to generate images for [OmniGen](https://github.com/VectorSpaceLab/OmniGen/).
The abstract from the paper is:
*The emergence of Large Language Models (LLMs) has unified language generation tasks and revolutionized human-machine interaction. However, in the realm of image generation, a unified model capable of handling various tasks within a single framework remains largely unexplored. In this work, we introduce OmniGen, a new diffusion model for unified image generation. OmniGen is characterized by the following features: 1) Unification: OmniGen not only demonstrates text-to-image generation capabilities but also inherently supports various downstream tasks, such as image editing, subject-driven generation, and visual conditional generation. 2) Simplicity: The architecture of OmniGen is highly simplified, eliminating the need for additional plugins. Moreover, compared to existing diffusion models, it is more user-friendly and can complete complex tasks end-to-end through instructions without the need for extra intermediate steps, greatly simplifying the image generation workflow. 3) Knowledge Transfer: Benefit from learning in a unified format, OmniGen effectively transfers knowledge across different tasks, manages unseen tasks and domains, and exhibits novel capabilities. We also explore the models reasoning capabilities and potential applications of the chain-of-thought mechanism. This work represents the first attempt at a general-purpose image generation model, and we will release our resources at https://github.com/VectorSpaceLab/OmniGen to foster future advancements.*
```python
import torch
from diffusers import OmniGenTransformer2DModel
transformer = OmniGenTransformer2DModel.from_pretrained("Shitao/OmniGen-v1-diffusers", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## OmniGenTransformer2DModel
[[autodoc]] OmniGenTransformer2DModel

View File

@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# WanTransformer3DModel
A Diffusion Transformer model for 3D video-like data was introduced in [Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
The model can be loaded with the following code snippet.
```python
from diffusers import WanTransformer3DModel
transformer = WanTransformer3DModel.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## WanTransformer3DModel
[[autodoc]] WanTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -29,43 +29,3 @@ Customized normalization layers for supporting various models in 🤗 Diffusers.
## AdaGroupNorm
[[autodoc]] models.normalization.AdaGroupNorm
## AdaLayerNormContinuous
[[autodoc]] models.normalization.AdaLayerNormContinuous
## RMSNorm
[[autodoc]] models.normalization.RMSNorm
## GlobalResponseNorm
[[autodoc]] models.normalization.GlobalResponseNorm
## LuminaLayerNormContinuous
[[autodoc]] models.normalization.LuminaLayerNormContinuous
## SD35AdaLayerNormZeroX
[[autodoc]] models.normalization.SD35AdaLayerNormZeroX
## AdaLayerNormZeroSingle
[[autodoc]] models.normalization.AdaLayerNormZeroSingle
## LuminaRMSNormZero
[[autodoc]] models.normalization.LuminaRMSNormZero
## LpNorm
[[autodoc]] models.normalization.LpNorm
## CogView3PlusAdaLayerNormZeroTextImage
[[autodoc]] models.normalization.CogView3PlusAdaLayerNormZeroTextImage
## CogVideoXLayerNormZero
[[autodoc]] models.normalization.CogVideoXLayerNormZero
## MochiRMSNormZero
[[autodoc]] models.transformers.transformer_mochi.MochiRMSNormZero
## MochiRMSNorm
[[autodoc]] models.normalization.MochiRMSNorm

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text-to-Video Generation with AnimateDiff
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
## Overview
[AnimateDiff: Animate Your Personalized Text-to-Image Diffusion Models without Specific Tuning](https://arxiv.org/abs/2307.04725) by Yuwei Guo, Ceyuan Yang, Anyi Rao, Yaohui Wang, Yu Qiao, Dahua Lin, Bo Dai.

View File

@@ -15,10 +15,6 @@
# CogVideoX
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer](https://arxiv.org/abs/2408.06072) from Tsinghua University & ZhipuAI, by Zhuoyi Yang, Jiayan Teng, Wendi Zheng, Ming Ding, Shiyu Huang, Jiazheng Xu, Yuanming Yang, Wenyi Hong, Xiaohan Zhang, Guanyu Feng, Da Yin, Xiaotao Gu, Yuxuan Zhang, Weihan Wang, Yean Cheng, Ting Liu, Bin Xu, Yuxiao Dong, Jie Tang.
The abstract from the paper is:

View File

@@ -1,34 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# CogView4
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [zRzRzRzRzRzRzR](https://github.com/zRzRzRzRzRzRzR). The original codebase can be found [here](https://huggingface.co/THUDM). The original weights can be found under [hf.co/THUDM](https://huggingface.co/THUDM).
## CogView4Pipeline
[[autodoc]] CogView4Pipeline
- all
- __call__
## CogView4PipelineOutput
[[autodoc]] pipelines.cogview4.pipeline_output.CogView4PipelineOutput

View File

@@ -15,10 +15,6 @@
# ConsisID
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Identity-Preserving Text-to-Video Generation by Frequency Decomposition](https://arxiv.org/abs/2411.17440) from Peking University & University of Rochester & etc, by Shenghai Yuan, Jinfa Huang, Xianyi He, Yunyang Ge, Yujun Shi, Liuhan Chen, Jiebo Luo, Li Yuan.
The abstract from the paper is:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# FluxControlInpaint
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
FluxControlInpaintPipeline is an implementation of Inpainting for Flux.1 Depth/Canny models. It is a pipeline that allows you to inpaint images using the Flux.1 Depth/Canny models. The pipeline takes an image and a mask as input and returns the inpainted image.
FLUX.1 Depth and Canny [dev] is a 12 billion parameter rectified flow transformer capable of generating an image based on a text description while following the structure of a given input image. **This is not a ControlNet model**.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet with Flux.1
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
FluxControlNetPipeline is an implementation of ControlNet for Flux.1.
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.

View File

@@ -1,36 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNet
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
The abstract from the paper is:
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
This pipeline was contributed by [ishan24](https://huggingface.co/ishan24). ❤️
The original codebase can be found at [NVlabs/Sana](https://github.com/NVlabs/Sana), and you can find official ControlNet checkpoints on [Efficient-Large-Model's](https://huggingface.co/Efficient-Large-Model) Hub profile.
## SanaControlNetPipeline
[[autodoc]] SanaControlNetPipeline
- all
- __call__
## SanaPipelineOutput
[[autodoc]] pipelines.sana.pipeline_output.SanaPipelineOutput

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet with Stable Diffusion 3
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
StableDiffusion3ControlNetPipeline is an implementation of ControlNet for Stable Diffusion 3.
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet with Stable Diffusion XL
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNetUnion
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNetUnionModel is an implementation of ControlNet for Stable Diffusion XL.
The ControlNet model was introduced in [ControlNetPlus](https://github.com/xinsir6/ControlNetPlus) by xinsir6. It supports multiple conditioning inputs without increasing computation.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet-XS
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
Like the original ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

View File

@@ -12,11 +12,6 @@ specific language governing permissions and limitations under the License.
# DeepFloyd IF
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
## Overview
DeepFloyd IF is a novel state-of-the-art open-source text-to-image model with a high degree of photorealism and language understanding.

View File

@@ -1,88 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# EasyAnimate
[EasyAnimate](https://github.com/aigc-apps/EasyAnimate) by Alibaba PAI.
The description from it's GitHub page:
*EasyAnimate is a pipeline based on the transformer architecture, designed for generating AI images and videos, and for training baseline models and Lora models for Diffusion Transformer. We support direct prediction from pre-trained EasyAnimate models, allowing for the generation of videos with various resolutions, approximately 6 seconds in length, at 8fps (EasyAnimateV5.1, 1 to 49 frames). Additionally, users can train their own baseline and Lora models for specific style transformations.*
This pipeline was contributed by [bubbliiiing](https://github.com/bubbliiiing). The original codebase can be found [here](https://huggingface.co/alibaba-pai). The original weights can be found under [hf.co/alibaba-pai](https://huggingface.co/alibaba-pai).
There are two official EasyAnimate checkpoints for text-to-video and video-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`alibaba-pai/EasyAnimateV5.1-12b-zh`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh) | torch.float16 |
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-InP`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-InP) | torch.float16 |
There is one official EasyAnimate checkpoints available for image-to-video and video-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-InP`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-InP) | torch.float16 |
There are two official EasyAnimate checkpoints available for control-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-Control`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-Control) | torch.float16 |
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-Control-Camera`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-Control-Camera) | torch.float16 |
For the EasyAnimateV5.1 series:
- Text-to-video (T2V) and Image-to-video (I2V) works for multiple resolutions. The width and height can vary from 256 to 1024.
- Both T2V and I2V models support generation with 1~49 frames and work best at this value. Exporting videos at 8 FPS is recommended.
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`EasyAnimatePipeline`] for inference with bitsandbytes.
```py
import torch
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, EasyAnimateTransformer3DModel, EasyAnimatePipeline
from diffusers.utils import export_to_video
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True)
transformer_8bit = EasyAnimateTransformer3DModel.from_pretrained(
"alibaba-pai/EasyAnimateV5.1-12b-zh",
subfolder="transformer",
quantization_config=quant_config,
torch_dtype=torch.float16,
)
pipeline = EasyAnimatePipeline.from_pretrained(
"alibaba-pai/EasyAnimateV5.1-12b-zh",
transformer=transformer_8bit,
torch_dtype=torch.float16,
device_map="balanced",
)
prompt = "A cat walks on the grass, realistic style."
negative_prompt = "bad detailed"
video = pipeline(prompt=prompt, negative_prompt=negative_prompt, num_frames=49, num_inference_steps=30).frames[0]
export_to_video(video, "cat.mp4", fps=8)
```
## EasyAnimatePipeline
[[autodoc]] EasyAnimatePipeline
- all
- __call__
## EasyAnimatePipelineOutput
[[autodoc]] pipelines.easyanimate.pipeline_output.EasyAnimatePipelineOutput

View File

@@ -12,11 +12,6 @@ specific language governing permissions and limitations under the License.
# Flux
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
Flux is a series of text-to-image generation models based on diffusion transformers. To know more about Flux, check out the original [blog post](https://blackforestlabs.ai/announcing-black-forest-labs/) by the creators of Flux, Black Forest Labs.
Original model checkpoints for Flux can be found [here](https://huggingface.co/black-forest-labs). Original inference code can be found [here](https://github.com/black-forest-labs/flux).
@@ -360,74 +355,8 @@ image.save('flux_ip_adapter_output.jpg')
<figcaption class="mt-2 text-sm text-center text-gray-500">IP-Adapter examples with prompt "wearing sunglasses"</figcaption>
</div>
## Optimize
Flux is a very large model and requires ~50GB of RAM/VRAM to load all the modeling components. Enable some of the optimizations below to lower the memory requirements.
### Group offloading
[Group offloading](../../optimization/memory#group-offloading) lowers VRAM usage by offloading groups of internal layers rather than the whole model or weights. You need to use [`~hooks.apply_group_offloading`] on all the model components of a pipeline. The `offload_type` parameter allows you to toggle between block and leaf-level offloading. Setting it to `leaf_level` offloads the lowest leaf-level parameters to the CPU instead of offloading at the module-level.
On CUDA devices that support asynchronous data streaming, set `use_stream=True` to overlap data transfer and computation to accelerate inference.
> [!TIP]
> It is possible to mix block and leaf-level offloading for different components in a pipeline.
```py
import torch
from diffusers import FluxPipeline
from diffusers.hooks import apply_group_offloading
model_id = "black-forest-labs/FLUX.1-dev"
dtype = torch.bfloat16
pipe = FluxPipeline.from_pretrained(
model_id,
torch_dtype=dtype,
)
apply_group_offloading(
pipe.transformer,
offload_type="leaf_level",
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
use_stream=True,
)
apply_group_offloading(
pipe.text_encoder,
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
offload_type="leaf_level",
use_stream=True,
)
apply_group_offloading(
pipe.text_encoder_2,
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
offload_type="leaf_level",
use_stream=True,
)
apply_group_offloading(
pipe.vae,
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
offload_type="leaf_level",
use_stream=True,
)
prompt="A cat wearing sunglasses and working as a lifeguard at pool."
generator = torch.Generator().manual_seed(181201)
image = pipe(
prompt,
width=576,
height=1024,
num_inference_steps=30,
generator=generator
).images[0]
image
```
### Running FP16 inference
## Running FP16 inference
Flux can generate high-quality images with FP16 (i.e. to accelerate inference on Turing/Volta GPUs) but produces different outputs compared to FP32/BF16. The issue is that some activations in the text encoders have to be clipped when running in FP16, which affects the overall image. Forcing text encoders to run with FP32 inference thus removes this output difference. See [here](https://github.com/huggingface/diffusers/pull/9097#issuecomment-2272292516) for details.
@@ -456,7 +385,7 @@ out = pipe(
out.save("image.png")
```
### Quantization
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.

View File

@@ -1,43 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# HiDreamImage
[HiDream-I1](https://huggingface.co/HiDream-ai) by HiDream.ai
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Available models
The following models are available for the [`HiDreamImagePipeline`](text-to-image) pipeline:
| Model name | Description |
|:---|:---|
| [`HiDream-ai/HiDream-I1-Full`](https://huggingface.co/HiDream-ai/HiDream-I1-Full) | - |
| [`HiDream-ai/HiDream-I1-Dev`](https://huggingface.co/HiDream-ai/HiDream-I1-Dev) | - |
| [`HiDream-ai/HiDream-I1-Fast`](https://huggingface.co/HiDream-ai/HiDream-I1-Fast) | - |
## HiDreamImagePipeline
[[autodoc]] HiDreamImagePipeline
- all
- __call__
## HiDreamImagePipelineOutput
[[autodoc]] pipelines.hidream_image.pipeline_output.HiDreamImagePipelineOutput

View File

@@ -14,10 +14,6 @@
# HunyuanVideo
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[HunyuanVideo](https://www.arxiv.org/abs/2412.03603) by Tencent.
*Recent advancements in video generation have significantly impacted daily life for both individuals and industries. However, the leading video generation models remain closed-source, resulting in a notable performance gap between industry capabilities and those available to the public. In this report, we introduce HunyuanVideo, an innovative open-source video foundation model that demonstrates performance in video generation comparable to, or even surpassing, that of leading closed-source models. HunyuanVideo encompasses a comprehensive framework that integrates several key elements, including data curation, advanced architectural design, progressive model scaling and training, and an efficient infrastructure tailored for large-scale model training and inference. As a result, we successfully trained a video generative model with over 13 billion parameters, making it the largest among all open-source models. We conducted extensive experiments and implemented a series of targeted designs to ensure high visual quality, motion dynamics, text-video alignment, and advanced filming techniques. According to evaluations by professionals, HunyuanVideo outperforms previous state-of-the-art models, including Runway Gen-3, Luma 1.6, and three top-performing Chinese video generative models. By releasing the code for the foundation model and its applications, we aim to bridge the gap between closed-source and open-source communities. This initiative will empower individuals within the community to experiment with their ideas, fostering a more dynamic and vibrant video generation ecosystem. The code is publicly available at [this https URL](https://github.com/tencent/HunyuanVideo).*
@@ -36,23 +32,6 @@ Recommendations for inference:
- For smaller resolution videos, try lower values of `shift` (between `2.0` to `5.0`) in the [Scheduler](https://huggingface.co/docs/diffusers/main/en/api/schedulers/flow_match_euler_discrete#diffusers.FlowMatchEulerDiscreteScheduler.shift). For larger resolution images, try higher values (between `7.0` and `12.0`). The default value is `7.0` for HunyuanVideo.
- For more information about supported resolutions and other details, please refer to the original repository [here](https://github.com/Tencent/HunyuanVideo/).
## Available models
The following models are available for the [`HunyuanVideoPipeline`](text-to-video) pipeline:
| Model name | Description |
|:---|:---|
| [`hunyuanvideo-community/HunyuanVideo`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo) | Official HunyuanVideo (guidance-distilled). Performs best at multiple resolutions and frames. Performs best with `guidance_scale=6.0`, `true_cfg_scale=1.0` and without a negative prompt. |
| [`https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-T2V`](https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-T2V) | Skywork's custom finetune of HunyuanVideo (de-distilled). Performs best with `97x544x960` resolution, `guidance_scale=1.0`, `true_cfg_scale=6.0` and a negative prompt. |
The following models are available for the image-to-video pipeline:
| Model name | Description |
|:---|:---|
| [`Skywork/SkyReels-V1-Hunyuan-I2V`](https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-I2V) | Skywork's custom finetune of HunyuanVideo (de-distilled). Performs best with `97x544x960` resolution. Performs best at `97x544x960` resolution, `guidance_scale=1.0`, `true_cfg_scale=6.0` and a negative prompt. |
| [`hunyuanvideo-community/HunyuanVideo-I2V-33ch`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo-I2V) | Tecent's official HunyuanVideo 33-channel I2V model. Performs best at resolutions of 480, 720, 960, 1280. A higher `shift` value when initializing the scheduler is recommended (good values are between 7 and 20). |
| [`hunyuanvideo-community/HunyuanVideo-I2V`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo-I2V) | Tecent's official HunyuanVideo 16-channel I2V model. Performs best at resolutions of 480, 720, 960, 1280. A higher `shift` value when initializing the scheduler is recommended (good values are between 7 and 20) |
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.

View File

@@ -9,10 +9,6 @@ specific language governing permissions and limitations under the License.
# Kandinsky 3
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Kandinsky 3 is created by [Vladimir Arkhipkin](https://github.com/oriBetelgeuse),[Anastasia Maltseva](https://github.com/NastyaMittseva),[Igor Pavlov](https://github.com/boomb0om),[Andrei Filatov](https://github.com/anvilarth),[Arseniy Shakhmatov](https://github.com/cene555),[Andrey Kuznetsov](https://github.com/kuznetsoffandrey),[Denis Dimitrov](https://github.com/denndimitrov), [Zein Shaheen](https://github.com/zeinsh)
The description from it's GitHub page:

View File

@@ -12,11 +12,6 @@ specific language governing permissions and limitations under the License.
# Kolors: Effective Training of Diffusion Model for Photorealistic Text-to-Image Synthesis
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/kolors/kolors_header_collage.png)
Kolors is a large-scale text-to-image generation model based on latent diffusion, developed by [the Kuaishou Kolors team](https://github.com/Kwai-Kolors/Kolors). Trained on billions of text-image pairs, Kolors exhibits significant advantages over both open-source and closed-source models in visual quality, complex semantic accuracy, and text rendering for both Chinese and English characters. Furthermore, Kolors supports both Chinese and English inputs, demonstrating strong performance in understanding and generating Chinese-specific content. For more details, please refer to this [technical report](https://github.com/Kwai-Kolors/Kolors/blob/master/imgs/Kolors_paper.pdf).

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Latent Consistency Models
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Latent Consistency Models (LCMs) were proposed in [Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference](https://huggingface.co/papers/2310.04378) by Simian Luo, Yiqin Tan, Longbo Huang, Jian Li, and Hang Zhao.
The abstract of the paper is as follows:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# LEDITS++
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
LEDITS++ was proposed in [LEDITS++: Limitless Image Editing using Text-to-Image Models](https://huggingface.co/papers/2311.16711) by Manuel Brack, Felix Friedrich, Katharina Kornmeier, Linoy Tsaban, Patrick Schramowski, Kristian Kersting, Apolinário Passos.
The abstract from the paper is:

View File

@@ -14,11 +14,6 @@
# LTX Video
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
[LTX Video](https://huggingface.co/Lightricks/LTX-Video) is the first DiT-based video generation model capable of generating high-quality videos in real-time. It produces 24 FPS videos at a 768x512 resolution faster than they can be watched. Trained on a large-scale dataset of diverse videos, the model generates high-resolution videos with realistic and varied content. We provide a model for both text-to-video as well as image + text-to-video usecases.
<Tip>
@@ -33,7 +28,6 @@ Available models:
|:-------------:|:-----------------:|
| [`LTX Video 0.9.0`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.safetensors) | `torch.bfloat16` |
| [`LTX Video 0.9.1`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.1.safetensors) | `torch.bfloat16` |
| [`LTX Video 0.9.5`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.5.safetensors) | `torch.bfloat16` |
Note: The recommended dtype is for the transformer component. The VAE and text encoders can be either `torch.float32`, `torch.bfloat16` or `torch.float16` but the recommended dtype is `torch.bfloat16` as used in the original repository.
@@ -198,12 +192,6 @@ export_to_video(video, "ship.mp4", fps=24)
- all
- __call__
## LTXConditionPipeline
[[autodoc]] LTXConditionPipeline
- all
- __call__
## LTXPipelineOutput
[[autodoc]] pipelines.ltx.pipeline_output.LTXPipelineOutput

View File

@@ -58,10 +58,10 @@ Use [`torch.compile`](https://huggingface.co/docs/diffusers/main/en/tutorials/fa
First, load the pipeline:
```python
from diffusers import LuminaPipeline
from diffusers import LuminaText2ImgPipeline
import torch
pipeline = LuminaPipeline.from_pretrained(
pipeline = LuminaText2ImgPipeline.from_pretrained(
"Alpha-VLLM/Lumina-Next-SFT-diffusers", torch_dtype=torch.bfloat16
).to("cuda")
```
@@ -86,11 +86,11 @@ image = pipeline(prompt="Upper body of a young woman in a Victorian-era outfit w
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`LuminaPipeline`] for inference with bitsandbytes.
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`LuminaText2ImgPipeline`] for inference with bitsandbytes.
```py
import torch
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, Transformer2DModel, LuminaPipeline
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, Transformer2DModel, LuminaText2ImgPipeline
from transformers import BitsAndBytesConfig as BitsAndBytesConfig, T5EncoderModel
quant_config = BitsAndBytesConfig(load_in_8bit=True)
@@ -109,7 +109,7 @@ transformer_8bit = Transformer2DModel.from_pretrained(
torch_dtype=torch.float16,
)
pipeline = LuminaPipeline.from_pretrained(
pipeline = LuminaText2ImgPipeline.from_pretrained(
"Alpha-VLLM/Lumina-Next-SFT-diffusers",
text_encoder=text_encoder_8bit,
transformer=transformer_8bit,
@@ -122,9 +122,9 @@ image = pipeline(prompt).images[0]
image.save("lumina.png")
```
## LuminaPipeline
## LuminaText2ImgPipeline
[[autodoc]] LuminaPipeline
[[autodoc]] LuminaText2ImgPipeline
- all
- __call__

View File

@@ -1,87 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# Lumina2
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Lumina Image 2.0: A Unified and Efficient Image Generative Model](https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0) is a 2 billion parameter flow-based diffusion transformer capable of generating diverse images from text descriptions.
The abstract from the paper is:
*We introduce Lumina-Image 2.0, an advanced text-to-image model that surpasses previous state-of-the-art methods across multiple benchmarks, while also shedding light on its potential to evolve into a generalist vision intelligence model. Lumina-Image 2.0 exhibits three key properties: (1) Unification it adopts a unified architecture that treats text and image tokens as a joint sequence, enabling natural cross-modal interactions and facilitating task expansion. Besides, since high-quality captioners can provide semantically better-aligned text-image training pairs, we introduce a unified captioning system, UniCaptioner, which generates comprehensive and precise captions for the model. This not only accelerates model convergence but also enhances prompt adherence, variable-length prompt handling, and task generalization via prompt templates. (2) Efficiency to improve the efficiency of the unified architecture, we develop a set of optimization techniques that improve semantic learning and fine-grained texture generation during training while incorporating inference-time acceleration strategies without compromising image quality. (3) Transparency we open-source all training details, code, and models to ensure full reproducibility, aiming to bridge the gap between well-resourced closed-source research teams and independent developers.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Using Single File loading with Lumina Image 2.0
Single file loading for Lumina Image 2.0 is available for the `Lumina2Transformer2DModel`
```python
import torch
from diffusers import Lumina2Transformer2DModel, Lumina2Pipeline
ckpt_path = "https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0/blob/main/consolidated.00-of-01.pth"
transformer = Lumina2Transformer2DModel.from_single_file(
ckpt_path, torch_dtype=torch.bfloat16
)
pipe = Lumina2Pipeline.from_pretrained(
"Alpha-VLLM/Lumina-Image-2.0", transformer=transformer, torch_dtype=torch.bfloat16
)
pipe.enable_model_cpu_offload()
image = pipe(
"a cat holding a sign that says hello",
generator=torch.Generator("cpu").manual_seed(0),
).images[0]
image.save("lumina-single-file.png")
```
## Using GGUF Quantized Checkpoints with Lumina Image 2.0
GGUF Quantized checkpoints for the `Lumina2Transformer2DModel` can be loaded via `from_single_file` with the `GGUFQuantizationConfig`
```python
from diffusers import Lumina2Transformer2DModel, Lumina2Pipeline, GGUFQuantizationConfig
ckpt_path = "https://huggingface.co/calcuis/lumina-gguf/blob/main/lumina2-q4_0.gguf"
transformer = Lumina2Transformer2DModel.from_single_file(
ckpt_path,
quantization_config=GGUFQuantizationConfig(compute_dtype=torch.bfloat16),
torch_dtype=torch.bfloat16,
)
pipe = Lumina2Pipeline.from_pretrained(
"Alpha-VLLM/Lumina-Image-2.0", transformer=transformer, torch_dtype=torch.bfloat16
)
pipe.enable_model_cpu_offload()
image = pipe(
"a cat holding a sign that says hello",
generator=torch.Generator("cpu").manual_seed(0),
).images[0]
image.save("lumina-gguf.png")
```
## Lumina2Pipeline
[[autodoc]] Lumina2Pipeline
- all
- __call__

View File

@@ -1,6 +1,4 @@
<!--
Copyright 2023-2025 Marigold Team, ETH Zürich. All rights reserved.
Copyright 2024-2025 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 Marigold authors and The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,120 +10,67 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Marigold Computer Vision
# Marigold Pipelines for Computer Vision Tasks
![marigold](https://marigoldmonodepth.github.io/images/teaser_collage_compressed.jpg)
Marigold was proposed in
[Repurposing Diffusion-Based Image Generators for Monocular Depth Estimation](https://huggingface.co/papers/2312.02145),
a CVPR 2024 Oral paper by
[Bingxin Ke](http://www.kebingxin.com/),
[Anton Obukhov](https://www.obukhov.ai/),
[Shengyu Huang](https://shengyuh.github.io/),
[Nando Metzger](https://nandometzger.github.io/),
[Rodrigo Caye Daudt](https://rcdaudt.github.io/), and
[Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
The core idea is to **repurpose the generative prior of Text-to-Image Latent Diffusion Models (LDMs) for traditional
computer vision tasks**.
This approach was explored by fine-tuning Stable Diffusion for **Monocular Depth Estimation**, as demonstrated in the
teaser above.
Marigold was proposed in [Repurposing Diffusion-Based Image Generators for Monocular Depth Estimation](https://huggingface.co/papers/2312.02145), a CVPR 2024 Oral paper by [Bingxin Ke](http://www.kebingxin.com/), [Anton Obukhov](https://www.obukhov.ai/), [Shengyu Huang](https://shengyuh.github.io/), [Nando Metzger](https://nandometzger.github.io/), [Rodrigo Caye Daudt](https://rcdaudt.github.io/), and [Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
The idea is to repurpose the rich generative prior of Text-to-Image Latent Diffusion Models (LDMs) for traditional computer vision tasks.
Initially, this idea was explored to fine-tune Stable Diffusion for Monocular Depth Estimation, as shown in the teaser above.
Later,
- [Tianfu Wang](https://tianfwang.github.io/) trained the first Latent Consistency Model (LCM) of Marigold, which unlocked fast single-step inference;
- [Kevin Qu](https://www.linkedin.com/in/kevin-qu-b3417621b/?locale=en_US) extended the approach to Surface Normals Estimation;
- [Anton Obukhov](https://www.obukhov.ai/) contributed the pipelines and documentation into diffusers (enabled and supported by [YiYi Xu](https://yiyixuxu.github.io/) and [Sayak Paul](https://sayak.dev/)).
Marigold was later extended in the follow-up paper,
[Marigold: Affordable Adaptation of Diffusion-Based Image Generators for Image Analysis](https://huggingface.co/papers/2312.02145),
authored by
[Bingxin Ke](http://www.kebingxin.com/),
[Kevin Qu](https://www.linkedin.com/in/kevin-qu-b3417621b/?locale=en_US),
[Tianfu Wang](https://tianfwang.github.io/),
[Nando Metzger](https://nandometzger.github.io/),
[Shengyu Huang](https://shengyuh.github.io/),
[Bo Li](https://www.linkedin.com/in/bobboli0202/),
[Anton Obukhov](https://www.obukhov.ai/), and
[Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
This work expanded Marigold to support new modalities such as **Surface Normals** and **Intrinsic Image Decomposition**
(IID), introduced a training protocol for **Latent Consistency Models** (LCM), and demonstrated **High-Resolution** (HR)
processing capability.
The abstract from the paper is:
<Tip>
The early Marigold models (`v1-0` and earlier) were optimized for best results with at least 10 inference steps.
LCM models were later developed to enable high-quality inference in just 1 to 4 steps.
Marigold models `v1-1` and later use the DDIM scheduler to achieve optimal
results in as few as 1 to 4 steps.
</Tip>
*Monocular depth estimation is a fundamental computer vision task. Recovering 3D depth from a single image is geometrically ill-posed and requires scene understanding, so it is not surprising that the rise of deep learning has led to a breakthrough. The impressive progress of monocular depth estimators has mirrored the growth in model capacity, from relatively modest CNNs to large Transformer architectures. Still, monocular depth estimators tend to struggle when presented with images with unfamiliar content and layout, since their knowledge of the visual world is restricted by the data seen during training, and challenged by zero-shot generalization to new domains. This motivates us to explore whether the extensive priors captured in recent generative diffusion models can enable better, more generalizable depth estimation. We introduce Marigold, a method for affine-invariant monocular depth estimation that is derived from Stable Diffusion and retains its rich prior knowledge. The estimator can be fine-tuned in a couple of days on a single GPU using only synthetic training data. It delivers state-of-the-art performance across a wide range of datasets, including over 20% performance gains in specific cases. Project page: https://marigoldmonodepth.github.io.*
## Available Pipelines
Each pipeline is tailored for a specific computer vision task, processing an input RGB image and generating a
corresponding prediction.
Currently, the following computer vision tasks are implemented:
Each pipeline supports one Computer Vision task, which takes an input RGB image as input and produces a *prediction* of the modality of interest, such as a depth map of the input image.
Currently, the following tasks are implemented:
| Pipeline | Predicted Modalities | Demos |
|---------------------------------------------------------------------------------------------------------------------------------------------|------------------------------------------------------------------------------------------------------------------|:--------------------------------------------------------------------------------------------------------------------------------------------------:|
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-lcm), [Slow Original Demo (DDIM)](https://huggingface.co/spaces/prs-eth/marigold) |
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-normals-lcm) |
| Pipeline | Recommended Model Checkpoints | Spaces (Interactive Apps) | Predicted Modalities |
|---------------------------------------------------------------------------------------------------------------------------------------------------|---------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|:------------------------------------------------------------------------------------:|------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [prs-eth/marigold-depth-v1-1](https://huggingface.co/prs-eth/marigold-depth-v1-1) | [Depth Estimation](https://huggingface.co/spaces/prs-eth/marigold) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) |
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [prs-eth/marigold-normals-v1-1](https://huggingface.co/prs-eth/marigold-normals-v1-1) | [Surface Normals Estimation](https://huggingface.co/spaces/prs-eth/marigold-normals) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) |
| [MarigoldIntrinsicsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_intrinsics.py) | [prs-eth/marigold-iid-appearance-v1-1](https://huggingface.co/prs-eth/marigold-iid-appearance-v1-1),<br>[prs-eth/marigold-iid-lighting-v1-1](https://huggingface.co/prs-eth/marigold-iid-lighting-v1-1) | [Intrinsic Image Decomposition](https://huggingface.co/spaces/prs-eth/marigold-iid) | [Albedo](https://en.wikipedia.org/wiki/Albedo), [Materials](https://www.n.aiq3d.com/wiki/roughnessmetalnessao-map), [Lighting](https://en.wikipedia.org/wiki/Diffuse_reflection) |
## Available Checkpoints
All original checkpoints are available under the [PRS-ETH](https://huggingface.co/prs-eth/) organization on Hugging Face.
They are designed for use with diffusers pipelines and the [original codebase](https://github.com/prs-eth/marigold), which can also be used to train
new model checkpoints.
The following is a summary of the recommended checkpoints, all of which produce reliable results with 1 to 4 steps.
| Checkpoint | Modality | Comment |
|-----------------------------------------------------------------------------------------------------|--------------|--------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
| [prs-eth/marigold-depth-v1-1](https://huggingface.co/prs-eth/marigold-depth-v1-1) | Depth | Affine-invariant depth prediction assigns each pixel a value between 0 (near plane) and 1 (far plane), with both planes determined by the model during inference. |
| [prs-eth/marigold-normals-v0-1](https://huggingface.co/prs-eth/marigold-normals-v0-1) | Normals | The surface normals predictions are unit-length 3D vectors in the screen space camera, with values in the range from -1 to 1. |
| [prs-eth/marigold-iid-appearance-v1-1](https://huggingface.co/prs-eth/marigold-iid-appearance-v1-1) | Intrinsics | InteriorVerse decomposition is comprised of Albedo and two BRDF material properties: Roughness and Metallicity. |
| [prs-eth/marigold-iid-lighting-v1-1](https://huggingface.co/prs-eth/marigold-iid-lighting-v1-1) | Intrinsics | HyperSim decomposition of an image &nbsp\\(I\\)&nbsp is comprised of Albedo &nbsp\\(A\\), Diffuse shading &nbsp\\(S\\), and Non-diffuse residual &nbsp\\(R\\): &nbsp\\(I = A*S+R\\). |
The original checkpoints can be found under the [PRS-ETH](https://huggingface.co/prs-eth/) Hugging Face organization.
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff
between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to
efficiently load the same components into multiple pipelines.
Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section
[here](../../using-diffusers/svd#reduce-memory-usage).
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines. Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section [here](../../using-diffusers/svd#reduce-memory-usage).
</Tip>
<Tip warning={true}>
Marigold pipelines were designed and tested with the scheduler embedded in the model checkpoint.
The optimal number of inference steps varies by scheduler, with no universal value that works best across all cases.
To accommodate this, the `num_inference_steps` parameter in the pipeline's `__call__` method defaults to `None` (see the
API reference).
Unless set explicitly, it inherits the value from the `default_denoising_steps` field in the checkpoint configuration
file (`model_index.json`).
This ensures high-quality predictions when invoking the pipeline with only the `image` argument.
Marigold pipelines were designed and tested only with `DDIMScheduler` and `LCMScheduler`.
Depending on the scheduler, the number of inference steps required to get reliable predictions varies, and there is no universal value that works best across schedulers.
Because of that, the default value of `num_inference_steps` in the `__call__` method of the pipeline is set to `None` (see the API reference).
Unless set explicitly, its value will be taken from the checkpoint configuration `model_index.json`.
This is done to ensure high-quality predictions when calling the pipeline with just the `image` argument.
</Tip>
See also Marigold [usage examples](../../using-diffusers/marigold_usage).
## Marigold Depth Prediction API
See also Marigold [usage examples](marigold_usage).
## MarigoldDepthPipeline
[[autodoc]] MarigoldDepthPipeline
- all
- __call__
## MarigoldNormalsPipeline
[[autodoc]] MarigoldNormalsPipeline
- all
- __call__
## MarigoldDepthOutput
[[autodoc]] pipelines.marigold.pipeline_marigold_depth.MarigoldDepthOutput
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_depth
## Marigold Normals Estimation API
[[autodoc]] MarigoldNormalsPipeline
- __call__
[[autodoc]] pipelines.marigold.pipeline_marigold_normals.MarigoldNormalsOutput
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_normals
## Marigold Intrinsic Image Decomposition API
[[autodoc]] MarigoldIntrinsicsPipeline
- __call__
[[autodoc]] pipelines.marigold.pipeline_marigold_intrinsics.MarigoldIntrinsicsOutput
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_intrinsics
## MarigoldNormalsOutput
[[autodoc]] pipelines.marigold.pipeline_marigold_normals.MarigoldNormalsOutput

View File

@@ -15,10 +15,6 @@
# Mochi 1 Preview
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
> [!TIP]
> Only a research preview of the model weights is available at the moment.

View File

@@ -1,80 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# OmniGen
[OmniGen: Unified Image Generation](https://arxiv.org/pdf/2409.11340) from BAAI, by Shitao Xiao, Yueze Wang, Junjie Zhou, Huaying Yuan, Xingrun Xing, Ruiran Yan, Chaofan Li, Shuting Wang, Tiejun Huang, Zheng Liu.
The abstract from the paper is:
*The emergence of Large Language Models (LLMs) has unified language generation tasks and revolutionized human-machine interaction. However, in the realm of image generation, a unified model capable of handling various tasks within a single framework remains largely unexplored. In this work, we introduce OmniGen, a new diffusion model for unified image generation. OmniGen is characterized by the following features: 1) Unification: OmniGen not only demonstrates text-to-image generation capabilities but also inherently supports various downstream tasks, such as image editing, subject-driven generation, and visual conditional generation. 2) Simplicity: The architecture of OmniGen is highly simplified, eliminating the need for additional plugins. Moreover, compared to existing diffusion models, it is more user-friendly and can complete complex tasks end-to-end through instructions without the need for extra intermediate steps, greatly simplifying the image generation workflow. 3) Knowledge Transfer: Benefit from learning in a unified format, OmniGen effectively transfers knowledge across different tasks, manages unseen tasks and domains, and exhibits novel capabilities. We also explore the models reasoning capabilities and potential applications of the chain-of-thought mechanism. This work represents the first attempt at a general-purpose image generation model, and we will release our resources at https://github.com/VectorSpaceLab/OmniGen to foster future advancements.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [staoxiao](https://github.com/staoxiao). The original codebase can be found [here](https://github.com/VectorSpaceLab/OmniGen). The original weights can be found under [hf.co/shitao](https://huggingface.co/Shitao/OmniGen-v1).
## Inference
First, load the pipeline:
```python
import torch
from diffusers import OmniGenPipeline
pipe = OmniGenPipeline.from_pretrained("Shitao/OmniGen-v1-diffusers", torch_dtype=torch.bfloat16)
pipe.to("cuda")
```
For text-to-image, pass a text prompt. By default, OmniGen generates a 1024x1024 image.
You can try setting the `height` and `width` parameters to generate images with different size.
```python
prompt = "Realistic photo. A young woman sits on a sofa, holding a book and facing the camera. She wears delicate silver hoop earrings adorned with tiny, sparkling diamonds that catch the light, with her long chestnut hair cascading over her shoulders. Her eyes are focused and gentle, framed by long, dark lashes. She is dressed in a cozy cream sweater, which complements her warm, inviting smile. Behind her, there is a table with a cup of water in a sleek, minimalist blue mug. The background is a serene indoor setting with soft natural light filtering through a window, adorned with tasteful art and flowers, creating a cozy and peaceful ambiance. 4K, HD."
image = pipe(
prompt=prompt,
height=1024,
width=1024,
guidance_scale=3,
generator=torch.Generator(device="cpu").manual_seed(111),
).images[0]
image.save("output.png")
```
OmniGen supports multimodal inputs.
When the input includes an image, you need to add a placeholder `<img><|image_1|></img>` in the text prompt to represent the image.
It is recommended to enable `use_input_image_size_as_output` to keep the edited image the same size as the original image.
```python
prompt="<img><|image_1|></img> Remove the woman's earrings. Replace the mug with a clear glass filled with sparkling iced cola."
input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png")]
image = pipe(
prompt=prompt,
input_images=input_images,
guidance_scale=2,
img_guidance_scale=1.6,
use_input_image_size_as_output=True,
generator=torch.Generator(device="cpu").manual_seed(222)).images[0]
image.save("output.png")
```
## OmniGenPipeline
[[autodoc]] OmniGenPipeline
- all
- __call__

View File

@@ -54,7 +54,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [DiT](dit) | text2image |
| [Flux](flux) | text2image |
| [Hunyuan-DiT](hunyuandit) | text2image |
| [I2VGen-XL](i2vgenxl) | image2video |
| [I2VGen-XL](i2vgenxl) | text2video |
| [InstructPix2Pix](pix2pix) | image editing |
| [Kandinsky 2.1](kandinsky) | text2image, image2image, inpainting, interpolation |
| [Kandinsky 2.2](kandinsky_v22) | text2image, image2image, inpainting |
@@ -65,7 +65,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Latte](latte) | text2image |
| [LEDITS++](ledits_pp) | image editing |
| [Lumina-T2X](lumina) | text2image |
| [Marigold](marigold) | depth-estimation, normals-estimation, intrinsic-decomposition |
| [Marigold](marigold) | depth |
| [MultiDiffusion](panorama) | text2image |
| [MusicLDM](musicldm) | text2audio |
| [PAG](pag) | text2image |

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Perturbed-Attention Guidance
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Perturbed-Attention Guidance (PAG)](https://ku-cvlab.github.io/Perturbed-Attention-Guidance/) is a new diffusion sampling guidance that improves sample quality across both unconditional and conditional settings, achieving this without requiring further training or the integration of external modules.
PAG was introduced in [Self-Rectifying Diffusion Sampling with Perturbed-Attention Guidance](https://huggingface.co/papers/2403.17377) by Donghoon Ahn, Hyoungwon Cho, Jaewon Min, Wooseok Jang, Jungwoo Kim, SeonHwa Kim, Hyun Hee Park, Kyong Hwan Jin and Seungryong Kim.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# MultiDiffusion
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation](https://huggingface.co/papers/2302.08113) is by Omer Bar-Tal, Lior Yariv, Yaron Lipman, and Tali Dekel.
The abstract from the paper is:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Image-to-Video Generation with PIA (Personalized Image Animator)
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
## Overview
[PIA: Your Personalized Image Animator via Plug-and-Play Modules in Text-to-Image Models](https://arxiv.org/abs/2312.13964) by Yiming Zhang, Zhening Xing, Yanhong Zeng, Youqing Fang, Kai Chen

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# InstructPix2Pix
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[InstructPix2Pix: Learning to Follow Image Editing Instructions](https://huggingface.co/papers/2211.09800) is by Tim Brooks, Aleksander Holynski and Alexei A. Efros.
The abstract from the paper is:

View File

@@ -14,11 +14,6 @@
# SanaPipeline
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
[SANA: Efficient High-Resolution Image Synthesis with Linear Diffusion Transformers](https://huggingface.co/papers/2410.10629) from NVIDIA and MIT HAN Lab, by Enze Xie, Junsong Chen, Junyu Chen, Han Cai, Haotian Tang, Yujun Lin, Zhekai Zhang, Muyang Li, Ligeng Zhu, Yao Lu, Song Han.
The abstract from the paper is:

View File

@@ -1,100 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# SANA-Sprint
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[SANA-Sprint: One-Step Diffusion with Continuous-Time Consistency Distillation](https://huggingface.co/papers/2503.09641) from NVIDIA, MIT HAN Lab, and Hugging Face by Junsong Chen, Shuchen Xue, Yuyang Zhao, Jincheng Yu, Sayak Paul, Junyu Chen, Han Cai, Enze Xie, Song Han
The abstract from the paper is:
*This paper presents SANA-Sprint, an efficient diffusion model for ultra-fast text-to-image (T2I) generation. SANA-Sprint is built on a pre-trained foundation model and augmented with hybrid distillation, dramatically reducing inference steps from 20 to 1-4. We introduce three key innovations: (1) We propose a training-free approach that transforms a pre-trained flow-matching model for continuous-time consistency distillation (sCM), eliminating costly training from scratch and achieving high training efficiency. Our hybrid distillation strategy combines sCM with latent adversarial distillation (LADD): sCM ensures alignment with the teacher model, while LADD enhances single-step generation fidelity. (2) SANA-Sprint is a unified step-adaptive model that achieves high-quality generation in 1-4 steps, eliminating step-specific training and improving efficiency. (3) We integrate ControlNet with SANA-Sprint for real-time interactive image generation, enabling instant visual feedback for user interaction. SANA-Sprint establishes a new Pareto frontier in speed-quality tradeoffs, achieving state-of-the-art performance with 7.59 FID and 0.74 GenEval in only 1 step — outperforming FLUX-schnell (7.94 FID / 0.71 GenEval) while being 10× faster (0.1s vs 1.1s on H100). It also achieves 0.1s (T2I) and 0.25s (ControlNet) latency for 1024×1024 images on H100, and 0.31s (T2I) on an RTX 4090, showcasing its exceptional efficiency and potential for AI-powered consumer applications (AIPC). Code and pre-trained models will be open-sourced.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [lawrence-cj](https://github.com/lawrence-cj), [shuchen Xue](https://github.com/scxue) and [Enze Xie](https://github.com/xieenze). The original codebase can be found [here](https://github.com/NVlabs/Sana). The original weights can be found under [hf.co/Efficient-Large-Model](https://huggingface.co/Efficient-Large-Model/).
Available models:
| Model | Recommended dtype |
|:-------------------------------------------------------------------------------------------------------------------------------------------:|:-----------------:|
| [`Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers) | `torch.bfloat16` |
| [`Efficient-Large-Model/Sana_Sprint_0.6B_1024px_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_Sprint_0.6B_1024px_diffusers) | `torch.bfloat16` |
Refer to [this](https://huggingface.co/collections/Efficient-Large-Model/sana-sprint-67d6810d65235085b3b17c76) collection for more information.
Note: The recommended dtype mentioned is for the transformer weights. The text encoder must stay in `torch.bfloat16` and VAE weights must stay in `torch.bfloat16` or `torch.float32` for the model to work correctly. Please refer to the inference example below to see how to load the model with the recommended dtype.
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`SanaSprintPipeline`] for inference with bitsandbytes.
```py
import torch
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, SanaTransformer2DModel, SanaSprintPipeline
from transformers import BitsAndBytesConfig as BitsAndBytesConfig, AutoModel
quant_config = BitsAndBytesConfig(load_in_8bit=True)
text_encoder_8bit = AutoModel.from_pretrained(
"Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers",
subfolder="text_encoder",
quantization_config=quant_config,
torch_dtype=torch.bfloat16,
)
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True)
transformer_8bit = SanaTransformer2DModel.from_pretrained(
"Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers",
subfolder="transformer",
quantization_config=quant_config,
torch_dtype=torch.bfloat16,
)
pipeline = SanaSprintPipeline.from_pretrained(
"Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers",
text_encoder=text_encoder_8bit,
transformer=transformer_8bit,
torch_dtype=torch.bfloat16,
device_map="balanced",
)
prompt = "a tiny astronaut hatching from an egg on the moon"
image = pipeline(prompt).images[0]
image.save("sana.png")
```
## Setting `max_timesteps`
Users can tweak the `max_timesteps` value for experimenting with the visual quality of the generated outputs. The default `max_timesteps` value was obtained with an inference-time search process. For more details about it, check out the paper.
## SanaSprintPipeline
[[autodoc]] SanaSprintPipeline
- all
- __call__
## SanaPipelineOutput
[[autodoc]] pipelines.sana.pipeline_output.SanaPipelineOutput

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Depth-to-image
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model can also infer depth based on an image using [MiDaS](https://github.com/isl-org/MiDaS). This allows you to pass a text prompt and an initial image to condition the generation of new images as well as a `depth_map` to preserve the image structure.
<Tip>

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Image-to-image
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model can also be applied to image-to-image generation by passing a text prompt and an initial image to condition the generation of new images.
The [`StableDiffusionImg2ImgPipeline`] uses the diffusion-denoising mechanism proposed in [SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equations](https://huggingface.co/papers/2108.01073) by Chenlin Meng, Yutong He, Yang Song, Jiaming Song, Jiajun Wu, Jun-Yan Zhu, Stefano Ermon.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Inpainting
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model can also be applied to inpainting which lets you edit specific parts of an image by providing a mask and a text prompt using Stable Diffusion.
## Tips

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text-to-(RGB, depth)
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
LDM3D was proposed in [LDM3D: Latent Diffusion Model for 3D](https://huggingface.co/papers/2305.10853) by Gabriela Ben Melech Stan, Diana Wofk, Scottie Fox, Alex Redden, Will Saxton, Jean Yu, Estelle Aflalo, Shao-Yen Tseng, Fabio Nonato, Matthias Muller, and Vasudev Lal. LDM3D generates an image and a depth map from a given text prompt unlike the existing text-to-image diffusion models such as [Stable Diffusion](./overview) which only generates an image. With almost the same number of parameters, LDM3D achieves to create a latent space that can compress both the RGB images and the depth maps.
Two checkpoints are available for use:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable Diffusion pipelines
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/). Latent diffusion applies the diffusion process over a lower dimensional latent space to reduce memory and compute complexity. This specific type of diffusion model was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
Stable Diffusion is trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs.

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@@ -12,11 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable Diffusion 3
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
Stable Diffusion 3 (SD3) was proposed in [Scaling Rectified Flow Transformers for High-Resolution Image Synthesis](https://arxiv.org/pdf/2403.03206.pdf) by Patrick Esser, Sumith Kulal, Andreas Blattmann, Rahim Entezari, Jonas Muller, Harry Saini, Yam Levi, Dominik Lorenz, Axel Sauer, Frederic Boesel, Dustin Podell, Tim Dockhorn, Zion English, Kyle Lacey, Alex Goodwin, Yannik Marek, and Robin Rombach.
The abstract from the paper is:
@@ -82,7 +77,7 @@ from diffusers import StableDiffusion3Pipeline
from transformers import SiglipVisionModel, SiglipImageProcessor
image_encoder_id = "google/siglip-so400m-patch14-384"
ip_adapter_id = "InstantX/SD3.5-Large-IP-Adapter"
ip_adapter_id = "guiyrt/InstantX-SD3.5-Large-IP-Adapter-diffusers"
feature_extractor = SiglipImageProcessor.from_pretrained(
image_encoder_id,

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@@ -12,11 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable Diffusion XL
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
Stable Diffusion XL (SDXL) was proposed in [SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis](https://huggingface.co/papers/2307.01952) by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, and Robin Rombach.
The abstract from the paper is:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text-to-image
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model was created by researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [Runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionPipeline`] is capable of generating photorealistic images given any text input. It's trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs. Latent diffusion is the research on top of which Stable Diffusion was built. It was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
The abstract from the paper is:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Super-resolution
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion upscaler diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), and [LAION](https://laion.ai/). It is used to enhance the resolution of input images by a factor of 4.
<Tip>

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable unCLIP
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Stable unCLIP checkpoints are finetuned from [Stable Diffusion 2.1](./stable_diffusion/stable_diffusion_2) checkpoints to condition on CLIP image embeddings.
Stable unCLIP still conditions on text embeddings. Given the two separate conditionings, stable unCLIP can be used
for text guided image variation. When combined with an unCLIP prior, it can also be used for full text to image generation.

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@@ -18,10 +18,6 @@ specific language governing permissions and limitations under the License.
# Text-to-video
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[ModelScope Text-to-Video Technical Report](https://arxiv.org/abs/2308.06571) is by Jiuniu Wang, Hangjie Yuan, Dayou Chen, Yingya Zhang, Xiang Wang, Shiwei Zhang.
The abstract from the paper is:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text2Video-Zero
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Text2Video-Zero: Text-to-Image Diffusion Models are Zero-Shot Video Generators](https://huggingface.co/papers/2303.13439) is by Levon Khachatryan, Andranik Movsisyan, Vahram Tadevosyan, Roberto Henschel, [Zhangyang Wang](https://www.ece.utexas.edu/people/faculty/atlas-wang), Shant Navasardyan, [Humphrey Shi](https://www.humphreyshi.com).
Text2Video-Zero enables zero-shot video generation using either:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# UniDiffuser
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The UniDiffuser model was proposed in [One Transformer Fits All Distributions in Multi-Modal Diffusion at Scale](https://huggingface.co/papers/2303.06555) by Fan Bao, Shen Nie, Kaiwen Xue, Chongxuan Li, Shi Pu, Yaole Wang, Gang Yue, Yue Cao, Hang Su, Jun Zhu.
The abstract from the paper is:

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@@ -1,465 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# Wan
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
<!-- TODO(aryan): update abstract once paper is out -->
## Generating Videos with Wan 2.1
We will first need to install some addtional dependencies.
```shell
pip install -u ftfy imageio-ffmpeg imageio
```
### Text to Video Generation
The following example requires 11GB VRAM to run and uses the smaller `Wan-AI/Wan2.1-T2V-1.3B-Diffusers` model. You can switch it out
for the larger `Wan2.1-I2V-14B-720P-Diffusers` or `Wan-AI/Wan2.1-I2V-14B-480P-Diffusers` if you have at least 35GB VRAM available.
```python
from diffusers import WanPipeline
from diffusers.utils import export_to_video
# Available models: Wan-AI/Wan2.1-I2V-14B-720P-Diffusers or Wan-AI/Wan2.1-I2V-14B-480P-Diffusers
model_id = "Wan-AI/Wan2.1-T2V-1.3B-Diffusers"
pipe = WanPipeline.from_pretrained(model_id, torch_dtype=torch.bfloat16)
pipe.enable_model_cpu_offload()
prompt = "A cat and a dog baking a cake together in a kitchen. The cat is carefully measuring flour, while the dog is stirring the batter with a wooden spoon. The kitchen is cozy, with sunlight streaming through the window."
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
frames = pipe(prompt=prompt, negative_prompt=negative_prompt, num_frames=num_frames).frames[0]
export_to_video(frames, "wan-t2v.mp4", fps=16)
```
<Tip>
You can improve the quality of the generated video by running the decoding step in full precision.
</Tip>
```python
from diffusers import WanPipeline, AutoencoderKLWan
from diffusers.utils import export_to_video
model_id = "Wan-AI/Wan2.1-T2V-1.3B-Diffusers"
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
pipe = WanPipeline.from_pretrained(model_id, vae=vae, torch_dtype=torch.bfloat16)
# replace this with pipe.to("cuda") if you have sufficient VRAM
pipe.enable_model_cpu_offload()
prompt = "A cat and a dog baking a cake together in a kitchen. The cat is carefully measuring flour, while the dog is stirring the batter with a wooden spoon. The kitchen is cozy, with sunlight streaming through the window."
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
frames = pipe(prompt=prompt, num_frames=num_frames).frames[0]
export_to_video(frames, "wan-t2v.mp4", fps=16)
```
### Image to Video Generation
The Image to Video pipeline requires loading the `AutoencoderKLWan` and the `CLIPVisionModel` components in full precision. The following example will need at least
35GB of VRAM to run.
```python
import torch
import numpy as np
from diffusers import AutoencoderKLWan, WanImageToVideoPipeline
from diffusers.utils import export_to_video, load_image
from transformers import CLIPVisionModel
# Available models: Wan-AI/Wan2.1-I2V-14B-480P-Diffusers, Wan-AI/Wan2.1-I2V-14B-720P-Diffusers
model_id = "Wan-AI/Wan2.1-I2V-14B-480P-Diffusers"
image_encoder = CLIPVisionModel.from_pretrained(
model_id, subfolder="image_encoder", torch_dtype=torch.float32
)
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
pipe = WanImageToVideoPipeline.from_pretrained(
model_id, vae=vae, image_encoder=image_encoder, torch_dtype=torch.bfloat16
)
# replace this with pipe.to("cuda") if you have sufficient VRAM
pipe.enable_model_cpu_offload()
image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"
)
max_area = 480 * 832
aspect_ratio = image.height / image.width
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
image = image.resize((width, height))
prompt = (
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
)
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
output = pipe(
image=image,
prompt=prompt,
negative_prompt=negative_prompt,
height=height,
width=width,
num_frames=num_frames,
guidance_scale=5.0,
).frames[0]
export_to_video(output, "wan-i2v.mp4", fps=16)
```
### Video to Video Generation
```python
import torch
from diffusers.utils import load_video, export_to_video
from diffusers import AutoencoderKLWan, WanVideoToVideoPipeline, UniPCMultistepScheduler
# Available models: Wan-AI/Wan2.1-T2V-14B-Diffusers, Wan-AI/Wan2.1-T2V-1.3B-Diffusers
model_id = "Wan-AI/Wan2.1-T2V-1.3B-Diffusers"
vae = AutoencoderKLWan.from_pretrained(
model_id, subfolder="vae", torch_dtype=torch.float32
)
pipe = WanVideoToVideoPipeline.from_pretrained(
model_id, vae=vae, torch_dtype=torch.bfloat16
)
flow_shift = 3.0 # 5.0 for 720P, 3.0 for 480P
pipe.scheduler = UniPCMultistepScheduler.from_config(
pipe.scheduler.config, flow_shift=flow_shift
)
# change to pipe.to("cuda") if you have sufficient VRAM
pipe.enable_model_cpu_offload()
prompt = "A robot standing on a mountain top. The sun is setting in the background"
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
video = load_video(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/hiker.mp4"
)
output = pipe(
video=video,
prompt=prompt,
negative_prompt=negative_prompt,
height=480,
width=512,
guidance_scale=7.0,
strength=0.7,
).frames[0]
export_to_video(output, "wan-v2v.mp4", fps=16)
```
## Memory Optimizations for Wan 2.1
Base inference with the large 14B Wan 2.1 models can take up to 35GB of VRAM when generating videos at 720p resolution. We'll outline a few memory optimizations we can apply to reduce the VRAM required to run the model.
We'll use `Wan-AI/Wan2.1-I2V-14B-720P-Diffusers` model in these examples to demonstrate the memory savings, but the techniques are applicable to all model checkpoints.
### Group Offloading the Transformer and UMT5 Text Encoder
Find more information about group offloading [here](../optimization/memory.md)
#### Block Level Group Offloading
We can reduce our VRAM requirements by applying group offloading to the larger model components of the pipeline; the `WanTransformer3DModel` and `UMT5EncoderModel`. Group offloading will break up the individual modules of a model and offload/onload them onto your GPU as needed during inference. In this example, we'll apply `block_level` offloading, which will group the modules in a model into blocks of size `num_blocks_per_group` and offload/onload them to GPU. Moving to between CPU and GPU does add latency to the inference process. You can trade off between latency and memory savings by increasing or decreasing the `num_blocks_per_group`.
The following example will now only require 14GB of VRAM to run, but will take approximately 30 minutes to generate a video.
```python
import torch
import numpy as np
from diffusers import AutoencoderKLWan, WanTransformer3DModel, WanImageToVideoPipeline
from diffusers.hooks.group_offloading import apply_group_offloading
from diffusers.utils import export_to_video, load_image
from transformers import UMT5EncoderModel, CLIPVisionModel
# Available models: Wan-AI/Wan2.1-I2V-14B-480P-Diffusers, Wan-AI/Wan2.1-I2V-14B-720P-Diffusers
model_id = "Wan-AI/Wan2.1-I2V-14B-720P-Diffusers"
image_encoder = CLIPVisionModel.from_pretrained(
model_id, subfolder="image_encoder", torch_dtype=torch.float32
)
text_encoder = UMT5EncoderModel.from_pretrained(model_id, subfolder="text_encoder", torch_dtype=torch.bfloat16)
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
transformer = WanTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
onload_device = torch.device("cuda")
offload_device = torch.device("cpu")
apply_group_offloading(text_encoder,
onload_device=onload_device,
offload_device=offload_device,
offload_type="block_level",
num_blocks_per_group=4
)
transformer.enable_group_offload(
onload_device=onload_device,
offload_device=offload_device,
offload_type="block_level",
num_blocks_per_group=4,
)
pipe = WanImageToVideoPipeline.from_pretrained(
model_id,
vae=vae,
transformer=transformer,
text_encoder=text_encoder,
image_encoder=image_encoder,
torch_dtype=torch.bfloat16
)
# Since we've offloaded the larger models alrady, we can move the rest of the model components to GPU
pipe.to("cuda")
image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"
)
max_area = 720 * 832
aspect_ratio = image.height / image.width
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
image = image.resize((width, height))
prompt = (
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
)
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
output = pipe(
image=image,
prompt=prompt,
negative_prompt=negative_prompt,
height=height,
width=width,
num_frames=num_frames,
guidance_scale=5.0,
).frames[0]
export_to_video(output, "wan-i2v.mp4", fps=16)
```
#### Block Level Group Offloading with CUDA Streams
We can speed up group offloading inference, by enabling the use of [CUDA streams](https://pytorch.org/docs/stable/generated/torch.cuda.Stream.html). However, using CUDA streams requires moving the model parameters into pinned memory. This allocation is handled by Pytorch under the hood, and can result in a significant spike in CPU RAM usage. Please consider this option if your CPU RAM is atleast 2X the size of the model you are group offloading.
In the following example we will use CUDA streams when group offloading the `WanTransformer3DModel`. When testing on an A100, this example will require 14GB of VRAM, 52GB of CPU RAM, but will generate a video in approximately 9 minutes.
```python
import torch
import numpy as np
from diffusers import AutoencoderKLWan, WanTransformer3DModel, WanImageToVideoPipeline
from diffusers.hooks.group_offloading import apply_group_offloading
from diffusers.utils import export_to_video, load_image
from transformers import UMT5EncoderModel, CLIPVisionModel
# Available models: Wan-AI/Wan2.1-I2V-14B-480P-Diffusers, Wan-AI/Wan2.1-I2V-14B-720P-Diffusers
model_id = "Wan-AI/Wan2.1-I2V-14B-720P-Diffusers"
image_encoder = CLIPVisionModel.from_pretrained(
model_id, subfolder="image_encoder", torch_dtype=torch.float32
)
text_encoder = UMT5EncoderModel.from_pretrained(model_id, subfolder="text_encoder", torch_dtype=torch.bfloat16)
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
transformer = WanTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
onload_device = torch.device("cuda")
offload_device = torch.device("cpu")
apply_group_offloading(text_encoder,
onload_device=onload_device,
offload_device=offload_device,
offload_type="block_level",
num_blocks_per_group=4
)
transformer.enable_group_offload(
onload_device=onload_device,
offload_device=offload_device,
offload_type="leaf_level",
use_stream=True
)
pipe = WanImageToVideoPipeline.from_pretrained(
model_id,
vae=vae,
transformer=transformer,
text_encoder=text_encoder,
image_encoder=image_encoder,
torch_dtype=torch.bfloat16
)
# Since we've offloaded the larger models alrady, we can move the rest of the model components to GPU
pipe.to("cuda")
image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"
)
max_area = 720 * 832
aspect_ratio = image.height / image.width
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
image = image.resize((width, height))
prompt = (
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
)
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
output = pipe(
image=image,
prompt=prompt,
negative_prompt=negative_prompt,
height=height,
width=width,
num_frames=num_frames,
guidance_scale=5.0,
).frames[0]
export_to_video(output, "wan-i2v.mp4", fps=16)
```
### Applying Layerwise Casting to the Transformer
Find more information about layerwise casting [here](../optimization/memory.md)
In this example, we will model offloading with layerwise casting. Layerwise casting will downcast each layer's weights to `torch.float8_e4m3fn`, temporarily upcast to `torch.bfloat16` during the forward pass of the layer, then revert to `torch.float8_e4m3fn` afterward. This approach reduces memory requirements by approximately 50% while introducing a minor quality reduction in the generated video due to the precision trade-off.
This example will require 20GB of VRAM.
```python
import torch
import numpy as np
from diffusers import AutoencoderKLWan, WanTransformer3DModel, WanImageToVideoPipeline
from diffusers.hooks.group_offloading import apply_group_offloading
from diffusers.utils import export_to_video, load_image
from transformers import UMT5EncoderModel, CLIPVisionModel
model_id = "Wan-AI/Wan2.1-I2V-14B-720P-Diffusers"
image_encoder = CLIPVisionModel.from_pretrained(
model_id, subfolder="image_encoder", torch_dtype=torch.float32
)
text_encoder = UMT5EncoderModel.from_pretrained(model_id, subfolder="text_encoder", torch_dtype=torch.bfloat16)
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
transformer = WanTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
transformer.enable_layerwise_casting(storage_dtype=torch.float8_e4m3fn, compute_dtype=torch.bfloat16)
pipe = WanImageToVideoPipeline.from_pretrained(
model_id,
vae=vae,
transformer=transformer,
text_encoder=text_encoder,
image_encoder=image_encoder,
torch_dtype=torch.bfloat16
)
pipe.enable_model_cpu_offload()
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg")
max_area = 720 * 832
aspect_ratio = image.height / image.width
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
image = image.resize((width, height))
prompt = (
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
)
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
output = pipe(
image=image,
prompt=prompt,
negative_prompt=negative_prompt,
height=height,
width=width,
num_frames=num_frames,
num_inference_steps=50,
guidance_scale=5.0,
).frames[0]
export_to_video(output, "wan-i2v.mp4", fps=16)
```
## Using a Custom Scheduler
Wan can be used with many different schedulers, each with their own benefits regarding speed and generation quality. By default, Wan uses the `UniPCMultistepScheduler(prediction_type="flow_prediction", use_flow_sigmas=True, flow_shift=3.0)` scheduler. You can use a different scheduler as follows:
```python
from diffusers import FlowMatchEulerDiscreteScheduler, UniPCMultistepScheduler, WanPipeline
scheduler_a = FlowMatchEulerDiscreteScheduler(shift=5.0)
scheduler_b = UniPCMultistepScheduler(prediction_type="flow_prediction", use_flow_sigmas=True, flow_shift=4.0)
pipe = WanPipeline.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", scheduler=<CUSTOM_SCHEDULER_HERE>)
# or,
pipe.scheduler = <CUSTOM_SCHEDULER_HERE>
```
## Using Single File Loading with Wan 2.1
The `WanTransformer3DModel` and `AutoencoderKLWan` models support loading checkpoints in their original format via the `from_single_file` loading
method.
```python
import torch
from diffusers import WanPipeline, WanTransformer3DModel
ckpt_path = "https://huggingface.co/Comfy-Org/Wan_2.1_ComfyUI_repackaged/blob/main/split_files/diffusion_models/wan2.1_t2v_1.3B_bf16.safetensors"
transformer = WanTransformer3DModel.from_single_file(ckpt_path, torch_dtype=torch.bfloat16)
pipe = WanPipeline.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", transformer=transformer)
```
## Recommendations for Inference
- Keep `AutencoderKLWan` in `torch.float32` for better decoding quality.
- `num_frames` should satisfy the following constraint: `(num_frames - 1) % 4 == 0`
- For smaller resolution videos, try lower values of `shift` (between `2.0` to `5.0`) in the [Scheduler](https://huggingface.co/docs/diffusers/main/en/api/schedulers/flow_match_euler_discrete#diffusers.FlowMatchEulerDiscreteScheduler.shift). For larger resolution videos, try higher values (between `7.0` and `12.0`). The default value is `3.0` for Wan.
## WanPipeline
[[autodoc]] WanPipeline
- all
- __call__
## WanImageToVideoPipeline
[[autodoc]] WanImageToVideoPipeline
- all
- __call__
## WanPipelineOutput
[[autodoc]] pipelines.wan.pipeline_output.WanPipelineOutput

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Würstchen
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
<img src="https://github.com/dome272/Wuerstchen/assets/61938694/0617c863-165a-43ee-9303-2a17299a0cf9">
[Wuerstchen: An Efficient Architecture for Large-Scale Text-to-Image Diffusion Models](https://huggingface.co/papers/2306.00637) is by Pablo Pernias, Dominic Rampas, Mats L. Richter and Christopher Pal and Marc Aubreville.

View File

@@ -31,11 +31,6 @@ Learn how to quantize models in the [Quantization](../quantization/overview) gui
## GGUFQuantizationConfig
[[autodoc]] GGUFQuantizationConfig
## QuantoConfig
[[autodoc]] QuantoConfig
## TorchAoConfig
[[autodoc]] TorchAoConfig

View File

@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# CogVideoXDDIMScheduler
`CogVideoXDDIMScheduler` is based on [Denoising Diffusion Implicit Models](https://huggingface.co/papers/2010.02502), specifically for CogVideoX models.
## CogVideoXDDIMScheduler
[[autodoc]] CogVideoXDDIMScheduler

View File

@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# CogVideoXDPMScheduler
`CogVideoXDPMScheduler` is based on [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095), specifically for CogVideoX models.
## CogVideoXDPMScheduler
[[autodoc]] CogVideoXDPMScheduler

View File

@@ -45,7 +45,3 @@ Utility and helper functions for working with 🤗 Diffusers.
## apply_layerwise_casting
[[autodoc]] hooks.layerwise_casting.apply_layerwise_casting
## apply_group_offloading
[[autodoc]] hooks.group_offloading.apply_group_offloading

View File

@@ -83,8 +83,4 @@ Happy exploring, and thank you for being part of the Diffusers community!
<td><a href="https://github.com/suzukimain/auto_diffusers"> Model Search </a></td>
<td>Search models on Civitai and Hugging Face</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/beinsezii/skrample"> Skrample </a></td>
<td>Fully modular scheduler functions with 1st class diffusers integration.</td>
</tr>
</table>

View File

@@ -16,11 +16,6 @@ specific language governing permissions and limitations under the License.
<img src="https://colab.research.google.com/assets/colab-badge.svg" alt="Open In Colab"/>
</a>
> [!TIP]
> This document has now grown outdated given the emergence of existing evaluation frameworks for diffusion models for image generation. Please check
> out works like [HEIM](https://crfm.stanford.edu/helm/heim/latest/), [T2I-Compbench](https://arxiv.org/abs/2307.06350),
> [GenEval](https://arxiv.org/abs/2310.11513).
Evaluation of generative models like [Stable Diffusion](https://huggingface.co/docs/diffusers/stable_diffusion) is subjective in nature. But as practitioners and researchers, we often have to make careful choices amongst many different possibilities. So, when working with different generative models (like GANs, Diffusion, etc.), how do we choose one over the other?
Qualitative evaluation of such models can be error-prone and might incorrectly influence a decision.

View File

@@ -1,9 +0,0 @@
# Hybrid Inference API Reference
## Remote Decode
[[autodoc]] utils.remote_utils.remote_decode
## Remote Encode
[[autodoc]] utils.remote_utils.remote_encode

View File

@@ -1,60 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Hybrid Inference
**Empowering local AI builders with Hybrid Inference**
> [!TIP]
> Hybrid Inference is an [experimental feature](https://huggingface.co/blog/remote_vae).
> Feedback can be provided [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
## Why use Hybrid Inference?
Hybrid Inference offers a fast and simple way to offload local generation requirements.
- 🚀 **Reduced Requirements:** Access powerful models without expensive hardware.
- 💎 **Without Compromise:** Achieve the highest quality without sacrificing performance.
- 💰 **Cost Effective:** It's free! 🤑
- 🎯 **Diverse Use Cases:** Fully compatible with Diffusers 🧨 and the wider community.
- 🔧 **Developer-Friendly:** Simple requests, fast responses.
---
## Available Models
* **VAE Decode 🖼️:** Quickly decode latent representations into high-quality images without compromising performance or workflow speed.
* **VAE Encode 🔢:** Efficiently encode images into latent representations for generation and training.
* **Text Encoders 📃 (coming soon):** Compute text embeddings for your prompts quickly and accurately, ensuring a smooth and high-quality workflow.
---
## Integrations
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.
## Changelog
- March 10 2025: Added VAE encode
- March 2 2025: Initial release with VAE decoding
## Contents
The documentation is organized into three sections:
* **VAE Decode** Learn the basics of how to use VAE Decode with Hybrid Inference.
* **VAE Encode** Learn the basics of how to use VAE Encode with Hybrid Inference.
* **API Reference** Dive into task-specific settings and parameters.

View File

@@ -1,345 +0,0 @@
# Getting Started: VAE Decode with Hybrid Inference
VAE decode is an essential component of diffusion models - turning latent representations into images or videos.
## Memory
These tables demonstrate the VRAM requirements for VAE decode with SD v1 and SD XL on different GPUs.
For the majority of these GPUs the memory usage % dictates other models (text encoders, UNet/Transformer) must be offloaded, or tiled decoding has to be used which increases time taken and impacts quality.
<details><summary>SD v1.5</summary>
| GPU | Resolution | Time (seconds) | Memory (%) | Tiled Time (secs) | Tiled Memory (%) |
| --- | --- | --- | --- | --- | --- |
| NVIDIA GeForce RTX 4090 | 512x512 | 0.031 | 5.60% | 0.031 (0%) | 5.60% |
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.148 | 20.00% | 0.301 (+103%) | 5.60% |
| NVIDIA GeForce RTX 4080 | 512x512 | 0.05 | 8.40% | 0.050 (0%) | 8.40% |
| NVIDIA GeForce RTX 4080 | 1024x1024 | 0.224 | 30.00% | 0.356 (+59%) | 8.40% |
| NVIDIA GeForce RTX 4070 Ti | 512x512 | 0.066 | 11.30% | 0.066 (0%) | 11.30% |
| NVIDIA GeForce RTX 4070 Ti | 1024x1024 | 0.284 | 40.50% | 0.454 (+60%) | 11.40% |
| NVIDIA GeForce RTX 3090 | 512x512 | 0.062 | 5.20% | 0.062 (0%) | 5.20% |
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.253 | 18.50% | 0.464 (+83%) | 5.20% |
| NVIDIA GeForce RTX 3080 | 512x512 | 0.07 | 12.80% | 0.070 (0%) | 12.80% |
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.286 | 45.30% | 0.466 (+63%) | 12.90% |
| NVIDIA GeForce RTX 3070 | 512x512 | 0.102 | 15.90% | 0.102 (0%) | 15.90% |
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.421 | 56.30% | 0.746 (+77%) | 16.00% |
</details>
<details><summary>SDXL</summary>
| GPU | Resolution | Time (seconds) | Memory Consumed (%) | Tiled Time (seconds) | Tiled Memory (%) |
| --- | --- | --- | --- | --- | --- |
| NVIDIA GeForce RTX 4090 | 512x512 | 0.057 | 10.00% | 0.057 (0%) | 10.00% |
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.256 | 35.50% | 0.257 (+0.4%) | 35.50% |
| NVIDIA GeForce RTX 4080 | 512x512 | 0.092 | 15.00% | 0.092 (0%) | 15.00% |
| NVIDIA GeForce RTX 4080 | 1024x1024 | 0.406 | 53.30% | 0.406 (0%) | 53.30% |
| NVIDIA GeForce RTX 4070 Ti | 512x512 | 0.121 | 20.20% | 0.120 (-0.8%) | 20.20% |
| NVIDIA GeForce RTX 4070 Ti | 1024x1024 | 0.519 | 72.00% | 0.519 (0%) | 72.00% |
| NVIDIA GeForce RTX 3090 | 512x512 | 0.107 | 10.50% | 0.107 (0%) | 10.50% |
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.459 | 38.00% | 0.460 (+0.2%) | 38.00% |
| NVIDIA GeForce RTX 3080 | 512x512 | 0.121 | 25.60% | 0.121 (0%) | 25.60% |
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.524 | 93.00% | 0.524 (0%) | 93.00% |
| NVIDIA GeForce RTX 3070 | 512x512 | 0.183 | 31.80% | 0.183 (0%) | 31.80% |
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.794 | 96.40% | 0.794 (0%) | 96.40% |
</details>
## Available VAEs
| | **Endpoint** | **Model** |
|:-:|:-----------:|:--------:|
| **Stable Diffusion v1** | [https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud](https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud) | [`stabilityai/sd-vae-ft-mse`](https://hf.co/stabilityai/sd-vae-ft-mse) |
| **Stable Diffusion XL** | [https://x2dmsqunjd6k9prw.us-east-1.aws.endpoints.huggingface.cloud](https://x2dmsqunjd6k9prw.us-east-1.aws.endpoints.huggingface.cloud) | [`madebyollin/sdxl-vae-fp16-fix`](https://hf.co/madebyollin/sdxl-vae-fp16-fix) |
| **Flux** | [https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud](https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud) | [`black-forest-labs/FLUX.1-schnell`](https://hf.co/black-forest-labs/FLUX.1-schnell) |
| **HunyuanVideo** | [https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud](https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud) | [`hunyuanvideo-community/HunyuanVideo`](https://hf.co/hunyuanvideo-community/HunyuanVideo) |
> [!TIP]
> Model support can be requested [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
## Code
> [!TIP]
> Install `diffusers` from `main` to run the code: `pip install git+https://github.com/huggingface/diffusers@main`
A helper method simplifies interacting with Hybrid Inference.
```python
from diffusers.utils.remote_utils import remote_decode
```
### Basic example
Here, we show how to use the remote VAE on random tensors.
<details><summary>Code</summary>
```python
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=torch.randn([1, 4, 64, 64], dtype=torch.float16),
scaling_factor=0.18215,
)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/output.png"/>
</figure>
Usage for Flux is slightly different. Flux latents are packed so we need to send the `height` and `width`.
<details><summary>Code</summary>
```python
image = remote_decode(
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=torch.randn([1, 4096, 64], dtype=torch.float16),
height=1024,
width=1024,
scaling_factor=0.3611,
shift_factor=0.1159,
)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/flux_random_latent.png"/>
</figure>
Finally, an example for HunyuanVideo.
<details><summary>Code</summary>
```python
video = remote_decode(
endpoint="https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=torch.randn([1, 16, 3, 40, 64], dtype=torch.float16),
output_type="mp4",
)
with open("video.mp4", "wb") as f:
f.write(video)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<video
alt="queue.mp4"
autoplay loop autobuffer muted playsinline
>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/video_1.mp4" type="video/mp4">
</video>
</figure>
### Generation
But we want to use the VAE on an actual pipeline to get an actual image, not random noise. The example below shows how to do it with SD v1.5.
<details><summary>Code</summary>
```python
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
variant="fp16",
vae=None,
).to("cuda")
prompt = "Strawberry ice cream, in a stylish modern glass, coconut, splashing milk cream and honey, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious"
latent = pipe(
prompt=prompt,
output_type="latent",
).images
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
scaling_factor=0.18215,
)
image.save("test.jpg")
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/test.jpg"/>
</figure>
Heres another example with Flux.
<details><summary>Code</summary>
```python
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained(
"black-forest-labs/FLUX.1-schnell",
torch_dtype=torch.bfloat16,
vae=None,
).to("cuda")
prompt = "Strawberry ice cream, in a stylish modern glass, coconut, splashing milk cream and honey, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious"
latent = pipe(
prompt=prompt,
guidance_scale=0.0,
num_inference_steps=4,
output_type="latent",
).images
image = remote_decode(
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
height=1024,
width=1024,
scaling_factor=0.3611,
shift_factor=0.1159,
)
image.save("test.jpg")
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/test_1.jpg"/>
</figure>
Heres an example with HunyuanVideo.
<details><summary>Code</summary>
```python
from diffusers import HunyuanVideoPipeline, HunyuanVideoTransformer3DModel
model_id = "hunyuanvideo-community/HunyuanVideo"
transformer = HunyuanVideoTransformer3DModel.from_pretrained(
model_id, subfolder="transformer", torch_dtype=torch.bfloat16
)
pipe = HunyuanVideoPipeline.from_pretrained(
model_id, transformer=transformer, vae=None, torch_dtype=torch.float16
).to("cuda")
latent = pipe(
prompt="A cat walks on the grass, realistic",
height=320,
width=512,
num_frames=61,
num_inference_steps=30,
output_type="latent",
).frames
video = remote_decode(
endpoint="https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
output_type="mp4",
)
if isinstance(video, bytes):
with open("video.mp4", "wb") as f:
f.write(video)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<video
alt="queue.mp4"
autoplay loop autobuffer muted playsinline
>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/video.mp4" type="video/mp4">
</video>
</figure>
### Queueing
One of the great benefits of using a remote VAE is that we can queue multiple generation requests. While the current latent is being processed for decoding, we can already queue another one. This helps improve concurrency.
<details><summary>Code</summary>
```python
import queue
import threading
from IPython.display import display
from diffusers import StableDiffusionPipeline
def decode_worker(q: queue.Queue):
while True:
item = q.get()
if item is None:
break
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=item,
scaling_factor=0.18215,
)
display(image)
q.task_done()
q = queue.Queue()
thread = threading.Thread(target=decode_worker, args=(q,), daemon=True)
thread.start()
def decode(latent: torch.Tensor):
q.put(latent)
prompts = [
"Blueberry ice cream, in a stylish modern glass , ice cubes, nuts, mint leaves, splashing milk cream, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious",
"Lemonade in a glass, mint leaves, in an aqua and white background, flowers, ice cubes, halo, fluid motion, dynamic movement, soft lighting, digital painting, rule of thirds composition, Art by Greg rutkowski, Coby whitmore",
"Comic book art, beautiful, vintage, pastel neon colors, extremely detailed pupils, delicate features, light on face, slight smile, Artgerm, Mary Blair, Edmund Dulac, long dark locks, bangs, glowing, fashionable style, fairytale ambience, hot pink.",
"Masterpiece, vanilla cone ice cream garnished with chocolate syrup, crushed nuts, choco flakes, in a brown background, gold, cinematic lighting, Art by WLOP",
"A bowl of milk, falling cornflakes, berries, blueberries, in a white background, soft lighting, intricate details, rule of thirds, octane render, volumetric lighting",
"Cold Coffee with cream, crushed almonds, in a glass, choco flakes, ice cubes, wet, in a wooden background, cinematic lighting, hyper realistic painting, art by Carne Griffiths, octane render, volumetric lighting, fluid motion, dynamic movement, muted colors,",
]
pipe = StableDiffusionPipeline.from_pretrained(
"Lykon/dreamshaper-8",
torch_dtype=torch.float16,
vae=None,
).to("cuda")
pipe.unet = pipe.unet.to(memory_format=torch.channels_last)
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
_ = pipe(
prompt=prompts[0],
output_type="latent",
)
for prompt in prompts:
latent = pipe(
prompt=prompt,
output_type="latent",
).images
decode(latent)
q.put(None)
thread.join()
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<video
alt="queue.mp4"
autoplay loop autobuffer muted playsinline
>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/queue.mp4" type="video/mp4">
</video>
</figure>
## Integrations
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.

View File

@@ -1,183 +0,0 @@
# Getting Started: VAE Encode with Hybrid Inference
VAE encode is used for training, image-to-image and image-to-video - turning into images or videos into latent representations.
## Memory
These tables demonstrate the VRAM requirements for VAE encode with SD v1 and SD XL on different GPUs.
For the majority of these GPUs the memory usage % dictates other models (text encoders, UNet/Transformer) must be offloaded, or tiled encoding has to be used which increases time taken and impacts quality.
<details><summary>SD v1.5</summary>
| GPU | Resolution | Time (seconds) | Memory (%) | Tiled Time (secs) | Tiled Memory (%) |
|:------------------------------|:-------------|-----------------:|-------------:|--------------------:|-------------------:|
| NVIDIA GeForce RTX 4090 | 512x512 | 0.015 | 3.51901 | 0.015 | 3.51901 |
| NVIDIA GeForce RTX 4090 | 256x256 | 0.004 | 1.3154 | 0.005 | 1.3154 |
| NVIDIA GeForce RTX 4090 | 2048x2048 | 0.402 | 47.1852 | 0.496 | 3.51901 |
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.078 | 12.2658 | 0.094 | 3.51901 |
| NVIDIA GeForce RTX 4080 SUPER | 512x512 | 0.023 | 5.30105 | 0.023 | 5.30105 |
| NVIDIA GeForce RTX 4080 SUPER | 256x256 | 0.006 | 1.98152 | 0.006 | 1.98152 |
| NVIDIA GeForce RTX 4080 SUPER | 2048x2048 | 0.574 | 71.08 | 0.656 | 5.30105 |
| NVIDIA GeForce RTX 4080 SUPER | 1024x1024 | 0.111 | 18.4772 | 0.14 | 5.30105 |
| NVIDIA GeForce RTX 3090 | 512x512 | 0.032 | 3.52782 | 0.032 | 3.52782 |
| NVIDIA GeForce RTX 3090 | 256x256 | 0.01 | 1.31869 | 0.009 | 1.31869 |
| NVIDIA GeForce RTX 3090 | 2048x2048 | 0.742 | 47.3033 | 0.954 | 3.52782 |
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.136 | 12.2965 | 0.207 | 3.52782 |
| NVIDIA GeForce RTX 3080 | 512x512 | 0.036 | 8.51761 | 0.036 | 8.51761 |
| NVIDIA GeForce RTX 3080 | 256x256 | 0.01 | 3.18387 | 0.01 | 3.18387 |
| NVIDIA GeForce RTX 3080 | 2048x2048 | 0.863 | 86.7424 | 1.191 | 8.51761 |
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.157 | 29.6888 | 0.227 | 8.51761 |
| NVIDIA GeForce RTX 3070 | 512x512 | 0.051 | 10.6941 | 0.051 | 10.6941 |
| NVIDIA GeForce RTX 3070 | 256x256 | 0.015 | 3.99743 | 0.015 | 3.99743 |
| NVIDIA GeForce RTX 3070 | 2048x2048 | 1.217 | 96.054 | 1.482 | 10.6941 |
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.223 | 37.2751 | 0.327 | 10.6941 |
</details>
<details><summary>SDXL</summary>
| GPU | Resolution | Time (seconds) | Memory Consumed (%) | Tiled Time (seconds) | Tiled Memory (%) |
|:------------------------------|:-------------|-----------------:|----------------------:|-----------------------:|-------------------:|
| NVIDIA GeForce RTX 4090 | 512x512 | 0.029 | 4.95707 | 0.029 | 4.95707 |
| NVIDIA GeForce RTX 4090 | 256x256 | 0.007 | 2.29666 | 0.007 | 2.29666 |
| NVIDIA GeForce RTX 4090 | 2048x2048 | 0.873 | 66.3452 | 0.863 | 15.5649 |
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.142 | 15.5479 | 0.143 | 15.5479 |
| NVIDIA GeForce RTX 4080 SUPER | 512x512 | 0.044 | 7.46735 | 0.044 | 7.46735 |
| NVIDIA GeForce RTX 4080 SUPER | 256x256 | 0.01 | 3.4597 | 0.01 | 3.4597 |
| NVIDIA GeForce RTX 4080 SUPER | 2048x2048 | 1.317 | 87.1615 | 1.291 | 23.447 |
| NVIDIA GeForce RTX 4080 SUPER | 1024x1024 | 0.213 | 23.4215 | 0.214 | 23.4215 |
| NVIDIA GeForce RTX 3090 | 512x512 | 0.058 | 5.65638 | 0.058 | 5.65638 |
| NVIDIA GeForce RTX 3090 | 256x256 | 0.016 | 2.45081 | 0.016 | 2.45081 |
| NVIDIA GeForce RTX 3090 | 2048x2048 | 1.755 | 77.8239 | 1.614 | 18.4193 |
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.265 | 18.4023 | 0.265 | 18.4023 |
| NVIDIA GeForce RTX 3080 | 512x512 | 0.064 | 13.6568 | 0.064 | 13.6568 |
| NVIDIA GeForce RTX 3080 | 256x256 | 0.018 | 5.91728 | 0.018 | 5.91728 |
| NVIDIA GeForce RTX 3080 | 2048x2048 | OOM | OOM | 1.866 | 44.4717 |
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.302 | 44.4308 | 0.302 | 44.4308 |
| NVIDIA GeForce RTX 3070 | 512x512 | 0.093 | 17.1465 | 0.093 | 17.1465 |
| NVIDIA GeForce RTX 3070 | 256x256 | 0.025 | 7.42931 | 0.026 | 7.42931 |
| NVIDIA GeForce RTX 3070 | 2048x2048 | OOM | OOM | 2.674 | 55.8355 |
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.443 | 55.7841 | 0.443 | 55.7841 |
</details>
## Available VAEs
| | **Endpoint** | **Model** |
|:-:|:-----------:|:--------:|
| **Stable Diffusion v1** | [https://qc6479g0aac6qwy9.us-east-1.aws.endpoints.huggingface.cloud](https://qc6479g0aac6qwy9.us-east-1.aws.endpoints.huggingface.cloud) | [`stabilityai/sd-vae-ft-mse`](https://hf.co/stabilityai/sd-vae-ft-mse) |
| **Stable Diffusion XL** | [https://xjqqhmyn62rog84g.us-east-1.aws.endpoints.huggingface.cloud](https://xjqqhmyn62rog84g.us-east-1.aws.endpoints.huggingface.cloud) | [`madebyollin/sdxl-vae-fp16-fix`](https://hf.co/madebyollin/sdxl-vae-fp16-fix) |
| **Flux** | [https://ptccx55jz97f9zgo.us-east-1.aws.endpoints.huggingface.cloud](https://ptccx55jz97f9zgo.us-east-1.aws.endpoints.huggingface.cloud) | [`black-forest-labs/FLUX.1-schnell`](https://hf.co/black-forest-labs/FLUX.1-schnell) |
> [!TIP]
> Model support can be requested [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
## Code
> [!TIP]
> Install `diffusers` from `main` to run the code: `pip install git+https://github.com/huggingface/diffusers@main`
A helper method simplifies interacting with Hybrid Inference.
```python
from diffusers.utils.remote_utils import remote_encode
```
### Basic example
Let's encode an image, then decode it to demonstrate.
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"/>
</figure>
<details><summary>Code</summary>
```python
from diffusers.utils import load_image
from diffusers.utils.remote_utils import remote_decode
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg?download=true")
latent = remote_encode(
endpoint="https://ptccx55jz97f9zgo.us-east-1.aws.endpoints.huggingface.cloud/",
scaling_factor=0.3611,
shift_factor=0.1159,
)
decoded = remote_decode(
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
scaling_factor=0.3611,
shift_factor=0.1159,
)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/decoded.png"/>
</figure>
### Generation
Now let's look at a generation example, we'll encode the image, generate then remotely decode too!
<details><summary>Code</summary>
```python
import torch
from diffusers import StableDiffusionImg2ImgPipeline
from diffusers.utils import load_image
from diffusers.utils.remote_utils import remote_decode, remote_encode
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
variant="fp16",
vae=None,
).to("cuda")
init_image = load_image(
"https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
)
init_image = init_image.resize((768, 512))
init_latent = remote_encode(
endpoint="https://qc6479g0aac6qwy9.us-east-1.aws.endpoints.huggingface.cloud/",
image=init_image,
scaling_factor=0.18215,
)
prompt = "A fantasy landscape, trending on artstation"
latent = pipe(
prompt=prompt,
image=init_latent,
strength=0.75,
output_type="latent",
).images
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
scaling_factor=0.18215,
)
image.save("fantasy_landscape.jpg")
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/fantasy_landscape.png"/>
</figure>
## Integrations
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.

View File

@@ -161,10 +161,10 @@ Your Python environment will find the `main` version of 🤗 Diffusers on the ne
Model weights and files are downloaded from the Hub to a cache which is usually your home directory. You can change the cache location by specifying the `HF_HOME` or `HUGGINFACE_HUB_CACHE` environment variables or configuring the `cache_dir` parameter in methods like [`~DiffusionPipeline.from_pretrained`].
Cached files allow you to run 🤗 Diffusers offline. To prevent 🤗 Diffusers from connecting to the internet, set the `HF_HUB_OFFLINE` environment variable to `1` and 🤗 Diffusers will only load previously downloaded files in the cache.
Cached files allow you to run 🤗 Diffusers offline. To prevent 🤗 Diffusers from connecting to the internet, set the `HF_HUB_OFFLINE` environment variable to `True` and 🤗 Diffusers will only load previously downloaded files in the cache.
```shell
export HF_HUB_OFFLINE=1
export HF_HUB_OFFLINE=True
```
For more details about managing and cleaning the cache, take a look at the [caching](https://huggingface.co/docs/huggingface_hub/guides/manage-cache) guide.
@@ -179,16 +179,14 @@ Telemetry is only sent when loading models and pipelines from the Hub,
and it is not collected if you're loading local files.
We understand that not everyone wants to share additional information,and we respect your privacy.
You can disable telemetry collection by setting the `HF_HUB_DISABLE_TELEMETRY` environment variable from your terminal:
You can disable telemetry collection by setting the `DISABLE_TELEMETRY` environment variable from your terminal:
On Linux/MacOS:
```bash
export HF_HUB_DISABLE_TELEMETRY=1
export DISABLE_TELEMETRY=YES
```
On Windows:
```bash
set HF_HUB_DISABLE_TELEMETRY=1
set DISABLE_TELEMETRY=YES
```

View File

@@ -158,62 +158,6 @@ In order to properly offload models after they're called, it is required to run
</Tip>
## Group offloading
Group offloading is the middle ground between sequential and model offloading. It works by offloading groups of internal layers (either `torch.nn.ModuleList` or `torch.nn.Sequential`), which uses less memory than model-level offloading. It is also faster than sequential-level offloading because the number of device synchronizations is reduced.
To enable group offloading, call the [`~ModelMixin.enable_group_offload`] method on the model if it is a Diffusers model implementation. For any other model implementation, use [`~hooks.group_offloading.apply_group_offloading`]:
```python
import torch
from diffusers import CogVideoXPipeline
from diffusers.hooks import apply_group_offloading
from diffusers.utils import export_to_video
# Load the pipeline
onload_device = torch.device("cuda")
offload_device = torch.device("cpu")
pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
# We can utilize the enable_group_offload method for Diffusers model implementations
pipe.transformer.enable_group_offload(onload_device=onload_device, offload_device=offload_device, offload_type="leaf_level", use_stream=True)
# Uncomment the following to also allow recording the current streams.
# pipe.transformer.enable_group_offload(onload_device=onload_device, offload_device=offload_device, offload_type="leaf_level", use_stream=True, record_stream=True)
# For any other model implementations, the apply_group_offloading function can be used
apply_group_offloading(pipe.text_encoder, onload_device=onload_device, offload_type="block_level", num_blocks_per_group=2)
apply_group_offloading(pipe.vae, onload_device=onload_device, offload_type="leaf_level")
prompt = (
"A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. "
"The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other "
"pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, "
"casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. "
"The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical "
"atmosphere of this unique musical performance."
)
video = pipe(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
# This utilized about 14.79 GB. It can be further reduced by using tiling and using leaf_level offloading throughout the pipeline.
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
export_to_video(video, "output.mp4", fps=8)
```
Group offloading (for CUDA devices with support for asynchronous data transfer streams) overlaps data transfer and computation to reduce the overall execution time compared to sequential offloading. This is enabled using layer prefetching with CUDA streams. The next layer to be executed is loaded onto the accelerator device while the current layer is being executed - this increases the memory requirements slightly. Group offloading also supports leaf-level offloading (equivalent to sequential CPU offloading) but can be made much faster when using streams.
<Tip>
- Group offloading may not work with all models out-of-the-box. If the forward implementations of the model contain weight-dependent device-casting of inputs, it may clash with the offloading mechanism's handling of device-casting.
- The `offload_type` parameter can be set to either `block_level` or `leaf_level`. `block_level` offloads groups of `torch::nn::ModuleList` or `torch::nn:Sequential` modules based on a configurable attribute `num_blocks_per_group`. For example, if you set `num_blocks_per_group=2` on a standard transformer model containing 40 layers, it will onload/offload 2 layers at a time for a total of 20 onload/offloads. This drastically reduces the VRAM requirements. `leaf_level` offloads individual layers at the lowest level, which is equivalent to sequential offloading. However, unlike sequential offloading, group offloading can be made much faster when using streams, with minimal compromise to end-to-end generation time.
- The `use_stream` parameter can be used with CUDA devices to enable prefetching layers for onload. It defaults to `False`. Layer prefetching allows overlapping computation and data transfer of model weights, which drastically reduces the overall execution time compared to other offloading methods. However, it can increase the CPU RAM usage significantly. Ensure that available CPU RAM that is at least twice the size of the model when setting `use_stream=True`. You can find more information about CUDA streams [here](https://pytorch.org/docs/stable/generated/torch.cuda.Stream.html)
- If specifying `use_stream=True` on VAEs with tiling enabled, make sure to do a dummy forward pass (possibly with dummy inputs) before the actual inference to avoid device-mismatch errors. This may not work on all implementations. Please open an issue if you encounter any problems.
- The parameter `low_cpu_mem_usage` can be set to `True` to reduce CPU memory usage when using streams for group offloading. This is useful when the CPU memory is the bottleneck, but it may counteract the benefits of using streams and increase the overall execution time. The CPU memory savings come from creating pinned-tensors on-the-fly instead of pre-pinning them. This parameter is better suited for using `leaf_level` offloading.
- When using `use_stream=True`, users can additionally specify `record_stream=True` to get better speedups at the expense of slightly increased memory usage. Refer to the [official PyTorch docs](https://pytorch.org/docs/stable/generated/torch.Tensor.record_stream.html) to know more about this.
For more information about available parameters and an explanation of how group offloading works, refer to [`~hooks.group_offloading.apply_group_offloading`].
</Tip>
## FP8 layerwise weight-casting
PyTorch supports `torch.float8_e4m3fn` and `torch.float8_e5m2` as weight storage dtypes, but they can't be used for computation in many different tensor operations due to unimplemented kernel support. However, you can use these dtypes to store model weights in fp8 precision and upcast them on-the-fly when the layers are used in the forward pass. This is known as layerwise weight-casting.
@@ -251,14 +195,6 @@ In the above example, layerwise casting is enabled on the transformer component
However, you gain more control and flexibility by directly utilizing the [`~hooks.layerwise_casting.apply_layerwise_casting`] function instead of [`~ModelMixin.enable_layerwise_casting`].
<Tip>
- Layerwise casting may not work with all models out-of-the-box. Sometimes, the forward implementations of the model might contain internal typecasting of weight values. Such implementations are not supported due to the currently simplistic implementation of layerwise casting, which assumes that the forward pass is independent of the weight precision and that the input dtypes are always in `compute_dtype`. An example of an incompatible implementation can be found [here](https://github.com/huggingface/transformers/blob/7f5077e53682ca855afc826162b204ebf809f1f9/src/transformers/models/t5/modeling_t5.py#L294-L299).
- Layerwise casting may fail on custom modeling implementations that make use of [PEFT](https://github.com/huggingface/peft) layers. Some minimal checks to handle this case is implemented but is not extensively tested or guaranteed to work in all cases.
- It can be also be applied partially to specific layers of a model. Partially applying layerwise casting can either be done manually by calling the `apply_layerwise_casting` function on specific internal modules, or by specifying the `skip_modules_pattern` and `skip_modules_classes` parameters for a root module. These parameters are particularly useful for layers such as normalization and modulation.
</Tip>
## Channels-last memory format
The channels-last memory format is an alternative way of ordering NCHW tensors in memory to preserve dimension ordering. Channels-last tensors are ordered in such a way that the channels become the densest dimension (storing images pixel-per-pixel). Since not all operators currently support the channels-last format, it may result in worst performance but you should still try and see if it works for your model.

View File

@@ -12,9 +12,6 @@ specific language governing permissions and limitations under the License.
# Metal Performance Shaders (MPS)
> [!TIP]
> Pipelines with a <img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22"> badge indicate a model can take advantage of the MPS backend on Apple silicon devices for faster inference. Feel free to open a [Pull Request](https://github.com/huggingface/diffusers/compare) to add this badge to pipelines that are missing it.
🤗 Diffusers is compatible with Apple silicon (M1/M2 chips) using the PyTorch [`mps`](https://pytorch.org/docs/stable/notes/mps.html) device, which uses the Metal framework to leverage the GPU on MacOS devices. You'll need to have:
- macOS computer with Apple silicon (M1/M2) hardware
@@ -40,7 +37,7 @@ image
<Tip warning={true}>
The PyTorch [mps](https://pytorch.org/docs/stable/notes/mps.html) backend does not support NDArray sizes greater than `2**32`. Please open an [Issue](https://github.com/huggingface/diffusers/issues/new/choose) if you encounter this problem so we can investigate.
Generating multiple prompts in a batch can [crash](https://github.com/huggingface/diffusers/issues/363) or fail to work reliably. We believe this is related to the [`mps`](https://github.com/pytorch/pytorch/issues/84039) backend in PyTorch. While this is being investigated, you should iterate instead of batching.
</Tip>
@@ -62,10 +59,6 @@ If you're using **PyTorch 1.13**, you need to "prime" the pipeline with an addit
## Troubleshoot
This section lists some common issues with using the `mps` backend and how to solve them.
### Attention slicing
M1/M2 performance is very sensitive to memory pressure. When this occurs, the system automatically swaps if it needs to which significantly degrades performance.
To prevent this from happening, we recommend *attention slicing* to reduce memory pressure during inference and prevent swapping. This is especially relevant if your computer has less than 64GB of system RAM, or if you generate images at non-standard resolutions larger than 512×512 pixels. Call the [`~DiffusionPipeline.enable_attention_slicing`] function on your pipeline:
@@ -79,7 +72,3 @@ pipeline.enable_attention_slicing()
```
Attention slicing performs the costly attention operation in multiple steps instead of all at once. It usually improves performance by ~20% in computers without universal memory, but we've observed *better performance* in most Apple silicon computers unless you have 64GB of RAM or more.
### Batch inference
Generating multiple prompts in a batch can crash or fail to work reliably. If this is the case, try iterating instead of batching.

View File

@@ -49,7 +49,7 @@ For Ada and higher-series GPUs. we recommend changing `torch_dtype` to `torch.bf
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
from transformers import BitsAndBytesConfig as TransformersBitsAndBytesConfig
from diffusers import AutoModel
from diffusers import FluxTransformer2DModel
from transformers import T5EncoderModel
quant_config = TransformersBitsAndBytesConfig(load_in_8bit=True,)
@@ -63,7 +63,7 @@ text_encoder_2_8bit = T5EncoderModel.from_pretrained(
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True,)
transformer_8bit = AutoModel.from_pretrained(
transformer_8bit = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/FLUX.1-dev",
subfolder="transformer",
quantization_config=quant_config,
@@ -74,7 +74,7 @@ transformer_8bit = AutoModel.from_pretrained(
By default, all the other modules such as `torch.nn.LayerNorm` are converted to `torch.float16`. You can change the data type of these modules with the `torch_dtype` parameter.
```diff
transformer_8bit = AutoModel.from_pretrained(
transformer_8bit = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/FLUX.1-dev",
subfolder="transformer",
quantization_config=quant_config,
@@ -133,7 +133,7 @@ For Ada and higher-series GPUs. we recommend changing `torch_dtype` to `torch.bf
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
from transformers import BitsAndBytesConfig as TransformersBitsAndBytesConfig
from diffusers import AutoModel
from diffusers import FluxTransformer2DModel
from transformers import T5EncoderModel
quant_config = TransformersBitsAndBytesConfig(load_in_4bit=True,)
@@ -147,7 +147,7 @@ text_encoder_2_4bit = T5EncoderModel.from_pretrained(
quant_config = DiffusersBitsAndBytesConfig(load_in_4bit=True,)
transformer_4bit = AutoModel.from_pretrained(
transformer_4bit = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/FLUX.1-dev",
subfolder="transformer",
quantization_config=quant_config,
@@ -158,7 +158,7 @@ transformer_4bit = AutoModel.from_pretrained(
By default, all the other modules such as `torch.nn.LayerNorm` are converted to `torch.float16`. You can change the data type of these modules with the `torch_dtype` parameter.
```diff
transformer_4bit = AutoModel.from_pretrained(
transformer_4bit = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/FLUX.1-dev",
subfolder="transformer",
quantization_config=quant_config,
@@ -217,11 +217,11 @@ print(model.get_memory_footprint())
Quantized models can be loaded from the [`~ModelMixin.from_pretrained`] method without needing to specify the `quantization_config` parameters:
```py
from diffusers import AutoModel, BitsAndBytesConfig
from diffusers import FluxTransformer2DModel, BitsAndBytesConfig
quantization_config = BitsAndBytesConfig(load_in_4bit=True)
model_4bit = AutoModel.from_pretrained(
model_4bit = FluxTransformer2DModel.from_pretrained(
"hf-internal-testing/flux.1-dev-nf4-pkg", subfolder="transformer"
)
```
@@ -243,13 +243,13 @@ An "outlier" is a hidden state value greater than a certain threshold, and these
To find the best threshold for your model, we recommend experimenting with the `llm_int8_threshold` parameter in [`BitsAndBytesConfig`]:
```py
from diffusers import AutoModel, BitsAndBytesConfig
from diffusers import FluxTransformer2DModel, BitsAndBytesConfig
quantization_config = BitsAndBytesConfig(
load_in_8bit=True, llm_int8_threshold=10,
)
model_8bit = AutoModel.from_pretrained(
model_8bit = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/FLUX.1-dev",
subfolder="transformer",
quantization_config=quantization_config,
@@ -305,7 +305,7 @@ NF4 is a 4-bit data type from the [QLoRA](https://hf.co/papers/2305.14314) paper
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
from transformers import BitsAndBytesConfig as TransformersBitsAndBytesConfig
from diffusers import AutoModel
from diffusers import FluxTransformer2DModel
from transformers import T5EncoderModel
quant_config = TransformersBitsAndBytesConfig(
@@ -325,7 +325,7 @@ quant_config = DiffusersBitsAndBytesConfig(
bnb_4bit_quant_type="nf4",
)
transformer_4bit = AutoModel.from_pretrained(
transformer_4bit = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/FLUX.1-dev",
subfolder="transformer",
quantization_config=quant_config,
@@ -343,7 +343,7 @@ Nested quantization is a technique that can save additional memory at no additio
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
from transformers import BitsAndBytesConfig as TransformersBitsAndBytesConfig
from diffusers import AutoModel
from diffusers import FluxTransformer2DModel
from transformers import T5EncoderModel
quant_config = TransformersBitsAndBytesConfig(
@@ -363,7 +363,7 @@ quant_config = DiffusersBitsAndBytesConfig(
bnb_4bit_use_double_quant=True,
)
transformer_4bit = AutoModel.from_pretrained(
transformer_4bit = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/FLUX.1-dev",
subfolder="transformer",
quantization_config=quant_config,
@@ -379,7 +379,7 @@ Once quantized, you can dequantize a model to its original precision, but this m
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig
from transformers import BitsAndBytesConfig as TransformersBitsAndBytesConfig
from diffusers import AutoModel
from diffusers import FluxTransformer2DModel
from transformers import T5EncoderModel
quant_config = TransformersBitsAndBytesConfig(
@@ -399,7 +399,7 @@ quant_config = DiffusersBitsAndBytesConfig(
bnb_4bit_use_double_quant=True,
)
transformer_4bit = AutoModel.from_pretrained(
transformer_4bit = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/FLUX.1-dev",
subfolder="transformer",
quantization_config=quant_config,

View File

@@ -36,6 +36,5 @@ Diffusers currently supports the following quantization methods.
- [BitsandBytes](./bitsandbytes)
- [TorchAO](./torchao)
- [GGUF](./gguf)
- [Quanto](./quanto.md)
[This resource](https://huggingface.co/docs/transformers/main/en/quantization/overview#when-to-use-what) provides a good overview of the pros and cons of different quantization techniques.

View File

@@ -1,148 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Quanto
[Quanto](https://github.com/huggingface/optimum-quanto) is a PyTorch quantization backend for [Optimum](https://huggingface.co/docs/optimum/en/index). It has been designed with versatility and simplicity in mind:
- All features are available in eager mode (works with non-traceable models)
- Supports quantization aware training
- Quantized models are compatible with `torch.compile`
- Quantized models are Device agnostic (e.g CUDA,XPU,MPS,CPU)
In order to use the Quanto backend, you will first need to install `optimum-quanto>=0.2.6` and `accelerate`
```shell
pip install optimum-quanto accelerate
```
Now you can quantize a model by passing the `QuantoConfig` object to the `from_pretrained()` method. Although the Quanto library does allow quantizing `nn.Conv2d` and `nn.LayerNorm` modules, currently, Diffusers only supports quantizing the weights in the `nn.Linear` layers of a model. The following snippet demonstrates how to apply `float8` quantization with Quanto.
```python
import torch
from diffusers import FluxTransformer2DModel, QuantoConfig
model_id = "black-forest-labs/FLUX.1-dev"
quantization_config = QuantoConfig(weights_dtype="float8")
transformer = FluxTransformer2DModel.from_pretrained(
model_id,
subfolder="transformer",
quantization_config=quantization_config,
torch_dtype=torch.bfloat16,
)
pipe = FluxPipeline.from_pretrained(model_id, transformer=transformer, torch_dtype=torch_dtype)
pipe.to("cuda")
prompt = "A cat holding a sign that says hello world"
image = pipe(
prompt, num_inference_steps=50, guidance_scale=4.5, max_sequence_length=512
).images[0]
image.save("output.png")
```
## Skipping Quantization on specific modules
It is possible to skip applying quantization on certain modules using the `modules_to_not_convert` argument in the `QuantoConfig`. Please ensure that the modules passed in to this argument match the keys of the modules in the `state_dict`
```python
import torch
from diffusers import FluxTransformer2DModel, QuantoConfig
model_id = "black-forest-labs/FLUX.1-dev"
quantization_config = QuantoConfig(weights_dtype="float8", modules_to_not_convert=["proj_out"])
transformer = FluxTransformer2DModel.from_pretrained(
model_id,
subfolder="transformer",
quantization_config=quantization_config,
torch_dtype=torch.bfloat16,
)
```
## Using `from_single_file` with the Quanto Backend
`QuantoConfig` is compatible with `~FromOriginalModelMixin.from_single_file`.
```python
import torch
from diffusers import FluxTransformer2DModel, QuantoConfig
ckpt_path = "https://huggingface.co/black-forest-labs/FLUX.1-dev/blob/main/flux1-dev.safetensors"
quantization_config = QuantoConfig(weights_dtype="float8")
transformer = FluxTransformer2DModel.from_single_file(ckpt_path, quantization_config=quantization_config, torch_dtype=torch.bfloat16)
```
## Saving Quantized models
Diffusers supports serializing Quanto models using the `~ModelMixin.save_pretrained` method.
The serialization and loading requirements are different for models quantized directly with the Quanto library and models quantized
with Diffusers using Quanto as the backend. It is currently not possible to load models quantized directly with Quanto into Diffusers using `~ModelMixin.from_pretrained`
```python
import torch
from diffusers import FluxTransformer2DModel, QuantoConfig
model_id = "black-forest-labs/FLUX.1-dev"
quantization_config = QuantoConfig(weights_dtype="float8")
transformer = FluxTransformer2DModel.from_pretrained(
model_id,
subfolder="transformer",
quantization_config=quantization_config,
torch_dtype=torch.bfloat16,
)
# save quantized model to reuse
transformer.save_pretrained("<your quantized model save path>")
# you can reload your quantized model with
model = FluxTransformer2DModel.from_pretrained("<your quantized model save path>")
```
## Using `torch.compile` with Quanto
Currently the Quanto backend supports `torch.compile` for the following quantization types:
- `int8` weights
```python
import torch
from diffusers import FluxPipeline, FluxTransformer2DModel, QuantoConfig
model_id = "black-forest-labs/FLUX.1-dev"
quantization_config = QuantoConfig(weights_dtype="int8")
transformer = FluxTransformer2DModel.from_pretrained(
model_id,
subfolder="transformer",
quantization_config=quantization_config,
torch_dtype=torch.bfloat16,
)
transformer = torch.compile(transformer, mode="max-autotune", fullgraph=True)
pipe = FluxPipeline.from_pretrained(
model_id, transformer=transformer, torch_dtype=torch_dtype
)
pipe.to("cuda")
images = pipe("A cat holding a sign that says hello").images[0]
images.save("flux-quanto-compile.png")
```
## Supported Quantization Types
### Weights
- float8
- int8
- int4
- int2

View File

@@ -26,13 +26,13 @@ The example below only quantizes the weights to int8.
```python
import torch
from diffusers import FluxPipeline, AutoModel, TorchAoConfig
from diffusers import FluxPipeline, FluxTransformer2DModel, TorchAoConfig
model_id = "black-forest-labs/FLUX.1-dev"
dtype = torch.bfloat16
quantization_config = TorchAoConfig("int8wo")
transformer = AutoModel.from_pretrained(
transformer = FluxTransformer2DModel.from_pretrained(
model_id,
subfolder="transformer",
quantization_config=quantization_config,
@@ -99,10 +99,10 @@ To serialize a quantized model in a given dtype, first load the model with the d
```python
import torch
from diffusers import AutoModel, TorchAoConfig
from diffusers import FluxTransformer2DModel, TorchAoConfig
quantization_config = TorchAoConfig("int8wo")
transformer = AutoModel.from_pretrained(
transformer = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/Flux.1-Dev",
subfolder="transformer",
quantization_config=quantization_config,
@@ -115,9 +115,9 @@ To load a serialized quantized model, use the [`~ModelMixin.from_pretrained`] me
```python
import torch
from diffusers import FluxPipeline, AutoModel
from diffusers import FluxPipeline, FluxTransformer2DModel
transformer = AutoModel.from_pretrained("/path/to/flux_int8wo", torch_dtype=torch.bfloat16, use_safetensors=False)
transformer = FluxTransformer2DModel.from_pretrained("/path/to/flux_int8wo", torch_dtype=torch.bfloat16, use_safetensors=False)
pipe = FluxPipeline.from_pretrained("black-forest-labs/Flux.1-Dev", transformer=transformer, torch_dtype=torch.bfloat16)
pipe.to("cuda")
@@ -126,15 +126,15 @@ image = pipe(prompt, num_inference_steps=30, guidance_scale=7.0).images[0]
image.save("output.png")
```
If you are using `torch<=2.6.0`, some quantization methods, such as `uint4wo`, cannot be loaded directly and may result in an `UnpicklingError` when trying to load the models, but work as expected when saving them. In order to work around this, one can load the state dict manually into the model. Note, however, that this requires using `weights_only=False` in `torch.load`, so it should be run only if the weights were obtained from a trustable source.
Some quantization methods, such as `uint4wo`, cannot be loaded directly and may result in an `UnpicklingError` when trying to load the models, but work as expected when saving them. In order to work around this, one can load the state dict manually into the model. Note, however, that this requires using `weights_only=False` in `torch.load`, so it should be run only if the weights were obtained from a trustable source.
```python
import torch
from accelerate import init_empty_weights
from diffusers import FluxPipeline, AutoModel, TorchAoConfig
from diffusers import FluxPipeline, FluxTransformer2DModel, TorchAoConfig
# Serialize the model
transformer = AutoModel.from_pretrained(
transformer = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/Flux.1-Dev",
subfolder="transformer",
quantization_config=TorchAoConfig("uint4wo"),
@@ -146,13 +146,10 @@ transformer.save_pretrained("/path/to/flux_uint4wo", safe_serialization=False, m
# Load the model
state_dict = torch.load("/path/to/flux_uint4wo/diffusion_pytorch_model.bin", weights_only=False, map_location="cpu")
with init_empty_weights():
transformer = AutoModel.from_config("/path/to/flux_uint4wo/config.json")
transformer = FluxTransformer2DModel.from_config("/path/to/flux_uint4wo/config.json")
transformer.load_state_dict(state_dict, strict=True, assign=True)
```
> [!TIP]
> The [`AutoModel`] API is supported for PyTorch >= 2.6 as shown in the examples below.
## Resources
- [TorchAO Quantization API](https://github.com/pytorch/ao/blob/main/torchao/quantization/README.md)

View File

@@ -163,9 +163,6 @@ Models are initiated with the [`~ModelMixin.from_pretrained`] method which also
>>> model = UNet2DModel.from_pretrained(repo_id, use_safetensors=True)
```
> [!TIP]
> Use the [`AutoModel`] API to automatically select a model class if you're unsure of which one to use.
To access the model parameters, call `model.config`:
```py

View File

@@ -31,10 +31,10 @@ To adapt your text-to-image model for inpainting, you'll need to change the numb
Initialize a [`UNet2DConditionModel`] with the pretrained text-to-image model weights, and change `in_channels` to 9. Changing the number of `in_channels` means you need to set `ignore_mismatched_sizes=True` and `low_cpu_mem_usage=False` to avoid a size mismatch error because the shape is different now.
```py
from diffusers import AutoModel
from diffusers import UNet2DConditionModel
model_id = "stable-diffusion-v1-5/stable-diffusion-v1-5"
unet = AutoModel.from_pretrained(
unet = UNet2DConditionModel.from_pretrained(
model_id,
subfolder="unet",
in_channels=9,

View File

@@ -339,10 +339,7 @@ import torch
from huggingface_hub.repocard import RepoCard
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16,
).to("cuda")
model_id = "sayakpaul/custom-diffusion-cat-wooden-pot"
pipeline = DiffusionPipeline.from_pretrained("sayakpaul/custom-diffusion-cat-wooden-pot", torch_dtype=torch.float16).to("cuda")
pipeline.unet.load_attn_procs(model_id, weight_name="pytorch_custom_diffusion_weights.bin")
pipeline.load_textual_inversion(model_id, weight_name="<new1>.bin")
pipeline.load_textual_inversion(model_id, weight_name="<new2>.bin")

View File

@@ -165,10 +165,10 @@ flush()
Load the diffusion transformer next which has 12.5B parameters. This time, set `device_map="auto"` to automatically distribute the model across two 16GB GPUs. The `auto` strategy is backed by [Accelerate](https://hf.co/docs/accelerate/index) and available as a part of the [Big Model Inference](https://hf.co/docs/accelerate/concept_guides/big_model_inference) feature. It starts by distributing a model across the fastest device first (GPU) before moving to slower devices like the CPU and hard drive if needed. The trade-off of storing model parameters on slower devices is slower inference latency.
```py
from diffusers import AutoModel
from diffusers import FluxTransformer2DModel
import torch
transformer = AutoModel.from_pretrained(
transformer = FluxTransformer2DModel.from_pretrained(
"black-forest-labs/FLUX.1-dev",
subfolder="transformer",
device_map="auto",

View File

@@ -32,9 +32,9 @@ The denoiser checkpoint can also have multiple shards and supports inference tha
For example, let's save a sharded checkpoint for the [SDXL UNet](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0/tree/main/unet):
```python
from diffusers import AutoModel
from diffusers import UNet2DConditionModel
unet = AutoModel.from_pretrained(
unet = UNet2DConditionModel.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", subfolder="unet"
)
unet.save_pretrained("sdxl-unet-sharded", max_shard_size="5GB")
@@ -43,10 +43,10 @@ unet.save_pretrained("sdxl-unet-sharded", max_shard_size="5GB")
The size of the fp32 variant of the SDXL UNet checkpoint is ~10.4GB. Set the `max_shard_size` parameter to 5GB to create 3 shards. After saving, you can load them in [`StableDiffusionXLPipeline`]:
```python
from diffusers import AutoModel, StableDiffusionXLPipeline
from diffusers import UNet2DConditionModel, StableDiffusionXLPipeline
import torch
unet = AutoModel.from_pretrained(
unet = UNet2DConditionModel.from_pretrained(
"sayakpaul/sdxl-unet-sharded", torch_dtype=torch.float16
)
pipeline = StableDiffusionXLPipeline.from_pretrained(

View File

@@ -221,7 +221,3 @@ pipe.delete_adapters("toy")
pipe.get_active_adapters()
["pixel"]
```
## PeftInputAutocastDisableHook
[[autodoc]] hooks.layerwise_casting.PeftInputAutocastDisableHook

View File

@@ -157,84 +157,6 @@ pipeline(
)
```
## IP Adapter Cutoff
IP Adapter is an image prompt adapter that can be used for diffusion models without any changes to the underlying model. We can use the IP Adapter Cutoff Callback to disable the IP Adapter after a certain number of steps. To set up the callback, you need to specify the number of denoising steps after which the callback comes into effect. You can do so by using either one of these two arguments:
- `cutoff_step_ratio`: Float number with the ratio of the steps.
- `cutoff_step_index`: Integer number with the exact number of the step.
We need to download the diffusion model and load the ip_adapter for it as follows:
```py
from diffusers import AutoPipelineForText2Image
from diffusers.utils import load_image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16).to("cuda")
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="sdxl_models", weight_name="ip-adapter_sdxl.bin")
pipeline.set_ip_adapter_scale(0.6)
```
The setup for the callback should look something like this:
```py
from diffusers import AutoPipelineForText2Image
from diffusers.callbacks import IPAdapterScaleCutoffCallback
from diffusers.utils import load_image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16
).to("cuda")
pipeline.load_ip_adapter(
"h94/IP-Adapter",
subfolder="sdxl_models",
weight_name="ip-adapter_sdxl.bin"
)
pipeline.set_ip_adapter_scale(0.6)
callback = IPAdapterScaleCutoffCallback(
cutoff_step_ratio=None,
cutoff_step_index=5
)
image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_adapter_diner.png"
)
generator = torch.Generator(device="cuda").manual_seed(2628670641)
images = pipeline(
prompt="a tiger sitting in a chair drinking orange juice",
ip_adapter_image=image,
negative_prompt="deformed, ugly, wrong proportion, low res, bad anatomy, worst quality, low quality",
generator=generator,
num_inference_steps=50,
callback_on_step_end=callback,
).images
images[0].save("custom_callback_img.png")
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/without_callback.png" alt="generated image of a tiger sitting in a chair drinking orange juice" />
<figcaption class="mt-2 text-center text-sm text-gray-500">without IPAdapterScaleCutoffCallback</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/with_callback2.png" alt="generated image of a tiger sitting in a chair drinking orange juice with ip adapter callback" />
<figcaption class="mt-2 text-center text-sm text-gray-500">with IPAdapterScaleCutoffCallback</figcaption>
</div>
</div>
## Display image after each generation step
> [!TIP]

Some files were not shown because too many files have changed in this diff Show More