Compare commits

..

1 Commits

Author SHA1 Message Date
sayakpaul
1f67e4e7f6 add a test for checking effective custom gc. 2025-01-29 11:04:12 +05:30
487 changed files with 4681 additions and 48069 deletions

View File

@@ -1,38 +0,0 @@
name: "\U0001F31F Remote VAE"
description: Feedback for remote VAE pilot
labels: [ "Remote VAE" ]
body:
- type: textarea
id: positive
validations:
required: true
attributes:
label: Did you like the remote VAE solution?
description: |
If you liked it, we would appreciate it if you could elaborate what you liked.
- type: textarea
id: feedback
validations:
required: true
attributes:
label: What can be improved about the current solution?
description: |
Let us know the things you would like to see improved. Note that we will work optimizing the solution once the pilot is over and we have usage.
- type: textarea
id: others
validations:
required: true
attributes:
label: What other VAEs you would like to see if the pilot goes well?
description: |
Provide a list of the VAEs you would like to see in the future if the pilot goes well.
- type: textarea
id: additional-info
attributes:
label: Notify the members of the team
description: |
Tag the following folks when submitting this feedback: @hlky @sayakpaul

View File

@@ -418,8 +418,6 @@ jobs:
test_location: "gguf"
- backend: "torchao"
test_location: "torchao"
- backend: "optimum_quanto"
test_location: "quanto"
runs-on:
group: aws-g6e-xlarge-plus
container:

View File

@@ -1,51 +0,0 @@
name: PR Style Bot
on:
issue_comment:
types: [created]
permissions:
contents: write
pull-requests: write
jobs:
style:
uses: huggingface/huggingface_hub/.github/workflows/style-bot-action.yml@main
with:
python_quality_dependencies: "[quality]"
pre_commit_script_name: "Download and Compare files from the main branch"
pre_commit_script: |
echo "Downloading the files from the main branch"
curl -o main_Makefile https://raw.githubusercontent.com/huggingface/diffusers/main/Makefile
curl -o main_setup.py https://raw.githubusercontent.com/huggingface/diffusers/refs/heads/main/setup.py
curl -o main_check_doc_toc.py https://raw.githubusercontent.com/huggingface/diffusers/refs/heads/main/utils/check_doc_toc.py
echo "Compare the files and raise error if needed"
diff_failed=0
if ! diff -q main_Makefile Makefile; then
echo "Error: The Makefile has changed. Please ensure it matches the main branch."
diff_failed=1
fi
if ! diff -q main_setup.py setup.py; then
echo "Error: The setup.py has changed. Please ensure it matches the main branch."
diff_failed=1
fi
if ! diff -q main_check_doc_toc.py utils/check_doc_toc.py; then
echo "Error: The utils/check_doc_toc.py has changed. Please ensure it matches the main branch."
diff_failed=1
fi
if [ $diff_failed -eq 1 ]; then
echo "❌ Error happened as we detected changes in the files that should not be changed ❌"
exit 1
fi
echo "No changes in the files. Proceeding..."
rm -rf main_Makefile main_setup.py main_check_doc_toc.py
style_command: "make style && make quality"
secrets:
bot_token: ${{ secrets.GITHUB_TOKEN }}

View File

@@ -2,7 +2,8 @@ name: Fast tests for PRs
on:
pull_request:
branches: [main]
branches:
- main
paths:
- "src/diffusers/**.py"
- "benchmarks/**.py"
@@ -63,7 +64,6 @@ jobs:
run: |
python utils/check_copies.py
python utils/check_dummies.py
python utils/check_support_list.py
make deps_table_check_updated
- name: Check if failure
if: ${{ failure() }}
@@ -120,8 +120,7 @@ jobs:
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git --no-deps
python -m uv pip install accelerate
- name: Environment
run: |

View File

@@ -1,250 +0,0 @@
name: Fast GPU Tests on PR
on:
pull_request:
branches: main
paths:
- "src/diffusers/models/modeling_utils.py"
- "src/diffusers/models/model_loading_utils.py"
- "src/diffusers/pipelines/pipeline_utils.py"
- "src/diffusers/pipeline_loading_utils.py"
- "src/diffusers/loaders/lora_base.py"
- "src/diffusers/loaders/lora_pipeline.py"
- "src/diffusers/loaders/peft.py"
- "tests/pipelines/test_pipelines_common.py"
- "tests/models/test_modeling_common.py"
workflow_dispatch:
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
env:
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
HF_HUB_ENABLE_HF_TRANSFER: 1
PYTEST_TIMEOUT: 600
PIPELINE_USAGE_CUTOFF: 1000000000 # set high cutoff so that only always-test pipelines run
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
- name: Environment
run: |
python utils/print_env.py
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
run: |
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
echo $matrix
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: test-pipelines.json
path: reports
torch_pipelines_cuda_tests:
name: Torch Pipelines CUDA Tests
needs: setup_torch_cuda_pipeline_matrix
strategy:
fail-fast: false
max-parallel: 8
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
- name: Environment
run: |
python utils/print_env.py
- name: Extract tests
id: extract_tests
run: |
pattern=$(python utils/extract_tests_from_mixin.py --type pipeline)
echo "$pattern" > /tmp/test_pattern.txt
echo "pattern_file=/tmp/test_pattern.txt" >> $GITHUB_OUTPUT
- name: PyTorch CUDA checkpoint tests on Ubuntu
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
if [ "${{ matrix.module }}" = "ip_adapters" ]; then
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
tests/pipelines/${{ matrix.module }}
else
pattern=$(cat ${{ steps.extract_tests.outputs.pattern_file }})
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx and $pattern" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
tests/pipelines/${{ matrix.module }}
fi
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: pipeline_${{ matrix.module }}_test_reports
path: reports
torch_cuda_tests:
name: Torch CUDA Tests
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
defaults:
run:
shell: bash
strategy:
fail-fast: false
max-parallel: 2
matrix:
module: [models, schedulers, lora, others]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
- name: Environment
run: |
python utils/print_env.py
- name: Extract tests
id: extract_tests
run: |
pattern=$(python utils/extract_tests_from_mixin.py --type ${{ matrix.module }})
echo "$pattern" > /tmp/test_pattern.txt
echo "pattern_file=/tmp/test_pattern.txt" >> $GITHUB_OUTPUT
- name: Run PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
pattern=$(cat ${{ steps.extract_tests.outputs.pattern_file }})
if [ -z "$pattern" ]; then
python -m pytest -n 1 -sv --max-worker-restart=0 --dist=loadfile -k "not Flax and not Onnx" tests/${{ matrix.module }} \
--make-reports=tests_torch_cuda_${{ matrix.module }}
else
python -m pytest -n 1 -sv --max-worker-restart=0 --dist=loadfile -k "not Flax and not Onnx and $pattern" tests/${{ matrix.module }} \
--make-reports=tests_torch_cuda_${{ matrix.module }}
fi
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_cuda_${{ matrix.module }}_stats.txt
cat reports/tests_torch_cuda_${{ matrix.module }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: torch_cuda_test_reports_${{ matrix.module }}
path: reports
run_examples_tests:
name: Examples PyTorch CUDA tests on Ubuntu
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test,training]
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install timm
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=examples_torch_cuda examples/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/examples_torch_cuda_stats.txt
cat reports/examples_torch_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: examples_test_reports
path: reports

View File

@@ -349,6 +349,7 @@ jobs:
container:
image: diffusers/diffusers-pytorch-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -358,6 +359,7 @@ jobs:
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"

View File

@@ -7,8 +7,8 @@ on:
default: 'diffusers/diffusers-pytorch-cuda'
description: 'Name of the Docker image'
required: true
pr_number:
description: 'PR number to test on'
branch:
description: 'PR Branch to test on'
required: true
test:
description: 'Tests to run (e.g.: `tests/models`).'
@@ -43,8 +43,8 @@ jobs:
exit 1
fi
if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines|lora) ]]; then
echo "Error: The input string must contain either 'models', 'pipelines', or 'lora' after 'tests/'."
if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines) ]]; then
echo "Error: The input string must contain either 'models' or 'pipelines' after 'tests/'."
exit 1
fi
@@ -53,13 +53,13 @@ jobs:
exit 1
fi
echo "$PY_TEST"
shell: bash -e {0}
- name: Checkout PR branch
uses: actions/checkout@v4
with:
ref: refs/pull/${{ inputs.pr_number }}/head
ref: ${{ github.event.inputs.branch }}
repository: ${{ github.event.pull_request.head.repo.full_name }}
- name: Install pytest
run: |

View File

@@ -13,6 +13,3 @@ jobs:
fetch-depth: 0
- name: Secret Scanning
uses: trufflesecurity/trufflehog@main
with:
extra_args: --results=verified,unknown

View File

@@ -76,14 +76,6 @@
- local: advanced_inference/outpaint
title: Outpainting
title: Advanced inference
- sections:
- local: hybrid_inference/overview
title: Overview
- local: hybrid_inference/vae_decode
title: VAE Decode
- local: hybrid_inference/api_reference
title: API Reference
title: Hybrid Inference
- sections:
- local: using-diffusers/cogvideox
title: CogVideoX
@@ -97,8 +89,6 @@
title: Kandinsky
- local: using-diffusers/ip_adapter
title: IP-Adapter
- local: using-diffusers/omnigen
title: OmniGen
- local: using-diffusers/pag
title: PAG
- local: using-diffusers/controlnet
@@ -173,8 +163,6 @@
title: gguf
- local: quantization/torchao
title: torchao
- local: quantization/quanto
title: quanto
title: Quantization Methods
- sections:
- local: optimization/fp16
@@ -288,12 +276,8 @@
title: ConsisIDTransformer3DModel
- local: api/models/cogview3plus_transformer2d
title: CogView3PlusTransformer2DModel
- local: api/models/cogview4_transformer2d
title: CogView4Transformer2DModel
- local: api/models/dit_transformer2d
title: DiTTransformer2DModel
- local: api/models/easyanimate_transformer3d
title: EasyAnimateTransformer3DModel
- local: api/models/flux_transformer
title: FluxTransformer2DModel
- local: api/models/hunyuan_transformer2d
@@ -304,14 +288,10 @@
title: LatteTransformer3DModel
- local: api/models/lumina_nextdit2d
title: LuminaNextDiT2DModel
- local: api/models/lumina2_transformer2d
title: Lumina2Transformer2DModel
- local: api/models/ltx_video_transformer3d
title: LTXVideoTransformer3DModel
- local: api/models/mochi_transformer3d
title: MochiTransformer3DModel
- local: api/models/omnigen_transformer
title: OmniGenTransformer2DModel
- local: api/models/pixart_transformer2d
title: PixArtTransformer2DModel
- local: api/models/prior_transformer
@@ -326,8 +306,6 @@
title: Transformer2DModel
- local: api/models/transformer_temporal
title: TransformerTemporalModel
- local: api/models/wan_transformer_3d
title: WanTransformer3DModel
title: Transformers
- sections:
- local: api/models/stable_cascade_unet
@@ -356,12 +334,8 @@
title: AutoencoderKLHunyuanVideo
- local: api/models/autoencoderkl_ltx_video
title: AutoencoderKLLTXVideo
- local: api/models/autoencoderkl_magvit
title: AutoencoderKLMagvit
- local: api/models/autoencoderkl_mochi
title: AutoencoderKLMochi
- local: api/models/autoencoder_kl_wan
title: AutoencoderKLWan
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_dc
@@ -402,8 +376,6 @@
title: CogVideoX
- local: api/pipelines/cogview3
title: CogView3
- local: api/pipelines/cogview4
title: CogView4
- local: api/pipelines/consisid
title: ConsisID
- local: api/pipelines/consistency_models
@@ -436,8 +408,6 @@
title: DiffEdit
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/easyanimate
title: EasyAnimate
- local: api/pipelines/flux
title: Flux
- local: api/pipelines/control_flux_inpaint
@@ -468,8 +438,6 @@
title: LEDITS++
- local: api/pipelines/ltx_video
title: LTXVideo
- local: api/pipelines/lumina2
title: Lumina 2.0
- local: api/pipelines/lumina
title: Lumina-T2X
- local: api/pipelines/marigold
@@ -480,8 +448,6 @@
title: MultiDiffusion
- local: api/pipelines/musicldm
title: MusicLDM
- local: api/pipelines/omnigen
title: OmniGen
- local: api/pipelines/pag
title: PAG
- local: api/pipelines/paint_by_example
@@ -554,8 +520,6 @@
title: UniDiffuser
- local: api/pipelines/value_guided_sampling
title: Value-guided sampling
- local: api/pipelines/wan
title: Wan
- local: api/pipelines/wuerstchen
title: Wuerstchen
title: Pipelines
@@ -565,10 +529,6 @@
title: Overview
- local: api/schedulers/cm_stochastic_iterative
title: CMStochasticIterativeScheduler
- local: api/schedulers/ddim_cogvideox
title: CogVideoXDDIMScheduler
- local: api/schedulers/multistep_dpm_solver_cogvideox
title: CogVideoXDPMScheduler
- local: api/schedulers/consistency_decoder
title: ConsistencyDecoderScheduler
- local: api/schedulers/cosine_dpm

View File

@@ -25,16 +25,3 @@ Customized activation functions for supporting various models in 🤗 Diffusers.
## ApproximateGELU
[[autodoc]] models.activations.ApproximateGELU
## SwiGLU
[[autodoc]] models.activations.SwiGLU
## FP32SiLU
[[autodoc]] models.activations.FP32SiLU
## LinearActivation
[[autodoc]] models.activations.LinearActivation

View File

@@ -147,20 +147,3 @@ An attention processor is a class for applying different types of attention mech
## XLAFlashAttnProcessor2_0
[[autodoc]] models.attention_processor.XLAFlashAttnProcessor2_0
## XFormersJointAttnProcessor
[[autodoc]] models.attention_processor.XFormersJointAttnProcessor
## IPAdapterXFormersAttnProcessor
[[autodoc]] models.attention_processor.IPAdapterXFormersAttnProcessor
## FluxIPAdapterJointAttnProcessor2_0
[[autodoc]] models.attention_processor.FluxIPAdapterJointAttnProcessor2_0
## XLAFluxFlashAttnProcessor2_0
[[autodoc]] models.attention_processor.XLAFluxFlashAttnProcessor2_0

View File

@@ -20,10 +20,6 @@ LoRA is a fast and lightweight training method that inserts and trains a signifi
- [`FluxLoraLoaderMixin`] provides similar functions for [Flux](https://huggingface.co/docs/diffusers/main/en/api/pipelines/flux).
- [`CogVideoXLoraLoaderMixin`] provides similar functions for [CogVideoX](https://huggingface.co/docs/diffusers/main/en/api/pipelines/cogvideox).
- [`Mochi1LoraLoaderMixin`] provides similar functions for [Mochi](https://huggingface.co/docs/diffusers/main/en/api/pipelines/mochi).
- [`LTXVideoLoraLoaderMixin`] provides similar functions for [LTX-Video](https://huggingface.co/docs/diffusers/main/en/api/pipelines/ltx_video).
- [`SanaLoraLoaderMixin`] provides similar functions for [Sana](https://huggingface.co/docs/diffusers/main/en/api/pipelines/sana).
- [`HunyuanVideoLoraLoaderMixin`] provides similar functions for [HunyuanVideo](https://huggingface.co/docs/diffusers/main/en/api/pipelines/hunyuan_video).
- [`Lumina2LoraLoaderMixin`] provides similar functions for [Lumina2](https://huggingface.co/docs/diffusers/main/en/api/pipelines/lumina2).
- [`AmusedLoraLoaderMixin`] is for the [`AmusedPipeline`].
- [`LoraBaseMixin`] provides a base class with several utility methods to fuse, unfuse, unload, LoRAs and more.
@@ -57,22 +53,6 @@ To learn more about how to load LoRA weights, see the [LoRA](../../using-diffuse
[[autodoc]] loaders.lora_pipeline.Mochi1LoraLoaderMixin
## LTXVideoLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.LTXVideoLoraLoaderMixin
## SanaLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.SanaLoraLoaderMixin
## HunyuanVideoLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.HunyuanVideoLoraLoaderMixin
## Lumina2LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.Lumina2LoraLoaderMixin
## AmusedLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.AmusedLoraLoaderMixin

View File

@@ -1,32 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLWan
The 3D variational autoencoder (VAE) model with KL loss used in [Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLWan
vae = AutoencoderKLWan.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", subfolder="vae", torch_dtype=torch.float32)
```
## AutoencoderKLWan
[[autodoc]] AutoencoderKLWan
- decode
- all
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

View File

@@ -1,37 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLMagvit
The 3D variational autoencoder (VAE) model with KL loss used in [EasyAnimate](https://github.com/aigc-apps/EasyAnimate) was introduced by Alibaba PAI.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLMagvit
vae = AutoencoderKLMagvit.from_pretrained("alibaba-pai/EasyAnimateV5.1-12b-zh", subfolder="vae", torch_dtype=torch.float16).to("cuda")
```
## AutoencoderKLMagvit
[[autodoc]] AutoencoderKLMagvit
- decode
- encode
- all
## AutoencoderKLOutput
[[autodoc]] models.autoencoders.autoencoder_kl.AutoencoderKLOutput
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

View File

@@ -1,30 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# CogView4Transformer2DModel
A Diffusion Transformer model for 2D data from [CogView4]()
The model can be loaded with the following code snippet.
```python
from diffusers import CogView4Transformer2DModel
transformer = CogView4Transformer2DModel.from_pretrained("THUDM/CogView4-6B", subfolder="transformer", torch_dtype=torch.bfloat16).to("cuda")
```
## CogView4Transformer2DModel
[[autodoc]] CogView4Transformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,30 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# EasyAnimateTransformer3DModel
A Diffusion Transformer model for 3D data from [EasyAnimate](https://github.com/aigc-apps/EasyAnimate) was introduced by Alibaba PAI.
The model can be loaded with the following code snippet.
```python
from diffusers import EasyAnimateTransformer3DModel
transformer = EasyAnimateTransformer3DModel.from_pretrained("alibaba-pai/EasyAnimateV5.1-12b-zh", subfolder="transformer", torch_dtype=torch.float16).to("cuda")
```
## EasyAnimateTransformer3DModel
[[autodoc]] EasyAnimateTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# Lumina2Transformer2DModel
A Diffusion Transformer model for 3D video-like data was introduced in [Lumina Image 2.0](https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0) by Alpha-VLLM.
The model can be loaded with the following code snippet.
```python
from diffusers import Lumina2Transformer2DModel
transformer = Lumina2Transformer2DModel.from_pretrained("Alpha-VLLM/Lumina-Image-2.0", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## Lumina2Transformer2DModel
[[autodoc]] Lumina2Transformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,30 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# OmniGenTransformer2DModel
A Transformer model that accepts multimodal instructions to generate images for [OmniGen](https://github.com/VectorSpaceLab/OmniGen/).
The abstract from the paper is:
*The emergence of Large Language Models (LLMs) has unified language generation tasks and revolutionized human-machine interaction. However, in the realm of image generation, a unified model capable of handling various tasks within a single framework remains largely unexplored. In this work, we introduce OmniGen, a new diffusion model for unified image generation. OmniGen is characterized by the following features: 1) Unification: OmniGen not only demonstrates text-to-image generation capabilities but also inherently supports various downstream tasks, such as image editing, subject-driven generation, and visual conditional generation. 2) Simplicity: The architecture of OmniGen is highly simplified, eliminating the need for additional plugins. Moreover, compared to existing diffusion models, it is more user-friendly and can complete complex tasks end-to-end through instructions without the need for extra intermediate steps, greatly simplifying the image generation workflow. 3) Knowledge Transfer: Benefit from learning in a unified format, OmniGen effectively transfers knowledge across different tasks, manages unseen tasks and domains, and exhibits novel capabilities. We also explore the models reasoning capabilities and potential applications of the chain-of-thought mechanism. This work represents the first attempt at a general-purpose image generation model, and we will release our resources at https://github.com/VectorSpaceLab/OmniGen to foster future advancements.*
```python
import torch
from diffusers import OmniGenTransformer2DModel
transformer = OmniGenTransformer2DModel.from_pretrained("Shitao/OmniGen-v1-diffusers", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## OmniGenTransformer2DModel
[[autodoc]] OmniGenTransformer2DModel

View File

@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# WanTransformer3DModel
A Diffusion Transformer model for 3D video-like data was introduced in [Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
The model can be loaded with the following code snippet.
```python
from diffusers import WanTransformer3DModel
transformer = WanTransformer3DModel.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## WanTransformer3DModel
[[autodoc]] WanTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -29,43 +29,3 @@ Customized normalization layers for supporting various models in 🤗 Diffusers.
## AdaGroupNorm
[[autodoc]] models.normalization.AdaGroupNorm
## AdaLayerNormContinuous
[[autodoc]] models.normalization.AdaLayerNormContinuous
## RMSNorm
[[autodoc]] models.normalization.RMSNorm
## GlobalResponseNorm
[[autodoc]] models.normalization.GlobalResponseNorm
## LuminaLayerNormContinuous
[[autodoc]] models.normalization.LuminaLayerNormContinuous
## SD35AdaLayerNormZeroX
[[autodoc]] models.normalization.SD35AdaLayerNormZeroX
## AdaLayerNormZeroSingle
[[autodoc]] models.normalization.AdaLayerNormZeroSingle
## LuminaRMSNormZero
[[autodoc]] models.normalization.LuminaRMSNormZero
## LpNorm
[[autodoc]] models.normalization.LpNorm
## CogView3PlusAdaLayerNormZeroTextImage
[[autodoc]] models.normalization.CogView3PlusAdaLayerNormZeroTextImage
## CogVideoXLayerNormZero
[[autodoc]] models.normalization.CogVideoXLayerNormZero
## MochiRMSNormZero
[[autodoc]] models.transformers.transformer_mochi.MochiRMSNormZero
## MochiRMSNorm
[[autodoc]] models.normalization.MochiRMSNorm

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text-to-Video Generation with AnimateDiff
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
## Overview
[AnimateDiff: Animate Your Personalized Text-to-Image Diffusion Models without Specific Tuning](https://arxiv.org/abs/2307.04725) by Yuwei Guo, Ceyuan Yang, Anyi Rao, Yaohui Wang, Yu Qiao, Dahua Lin, Bo Dai.

View File

@@ -15,10 +15,6 @@
# CogVideoX
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer](https://arxiv.org/abs/2408.06072) from Tsinghua University & ZhipuAI, by Zhuoyi Yang, Jiayan Teng, Wendi Zheng, Ming Ding, Shiyu Huang, Jiazheng Xu, Yuanming Yang, Wenyi Hong, Xiaohan Zhang, Guanyu Feng, Da Yin, Xiaotao Gu, Yuxuan Zhang, Weihan Wang, Yean Cheng, Ting Liu, Bin Xu, Yuxiao Dong, Jie Tang.
The abstract from the paper is:

View File

@@ -1,34 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# CogView4
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [zRzRzRzRzRzRzR](https://github.com/zRzRzRzRzRzRzR). The original codebase can be found [here](https://huggingface.co/THUDM). The original weights can be found under [hf.co/THUDM](https://huggingface.co/THUDM).
## CogView4Pipeline
[[autodoc]] CogView4Pipeline
- all
- __call__
## CogView4PipelineOutput
[[autodoc]] pipelines.cogview4.pipeline_output.CogView4PipelineOutput

View File

@@ -15,10 +15,6 @@
# ConsisID
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Identity-Preserving Text-to-Video Generation by Frequency Decomposition](https://arxiv.org/abs/2411.17440) from Peking University & University of Rochester & etc, by Shenghai Yuan, Jinfa Huang, Xianyi He, Yunyang Ge, Yujun Shi, Liuhan Chen, Jiebo Luo, Li Yuan.
The abstract from the paper is:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# FluxControlInpaint
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
FluxControlInpaintPipeline is an implementation of Inpainting for Flux.1 Depth/Canny models. It is a pipeline that allows you to inpaint images using the Flux.1 Depth/Canny models. The pipeline takes an image and a mask as input and returns the inpainted image.
FLUX.1 Depth and Canny [dev] is a 12 billion parameter rectified flow transformer capable of generating an image based on a text description while following the structure of a given input image. **This is not a ControlNet model**.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet with Flux.1
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
FluxControlNetPipeline is an implementation of ControlNet for Flux.1.
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet with Stable Diffusion 3
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
StableDiffusion3ControlNetPipeline is an implementation of ControlNet for Stable Diffusion 3.
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet with Stable Diffusion XL
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNetUnion
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNetUnionModel is an implementation of ControlNet for Stable Diffusion XL.
The ControlNet model was introduced in [ControlNetPlus](https://github.com/xinsir6/ControlNetPlus) by xinsir6. It supports multiple conditioning inputs without increasing computation.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet-XS
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
Like the original ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# DeepFloyd IF
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
## Overview
DeepFloyd IF is a novel state-of-the-art open-source text-to-image model with a high degree of photorealism and language understanding.

View File

@@ -1,88 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# EasyAnimate
[EasyAnimate](https://github.com/aigc-apps/EasyAnimate) by Alibaba PAI.
The description from it's GitHub page:
*EasyAnimate is a pipeline based on the transformer architecture, designed for generating AI images and videos, and for training baseline models and Lora models for Diffusion Transformer. We support direct prediction from pre-trained EasyAnimate models, allowing for the generation of videos with various resolutions, approximately 6 seconds in length, at 8fps (EasyAnimateV5.1, 1 to 49 frames). Additionally, users can train their own baseline and Lora models for specific style transformations.*
This pipeline was contributed by [bubbliiiing](https://github.com/bubbliiiing). The original codebase can be found [here](https://huggingface.co/alibaba-pai). The original weights can be found under [hf.co/alibaba-pai](https://huggingface.co/alibaba-pai).
There are two official EasyAnimate checkpoints for text-to-video and video-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`alibaba-pai/EasyAnimateV5.1-12b-zh`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh) | torch.float16 |
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-InP`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-InP) | torch.float16 |
There is one official EasyAnimate checkpoints available for image-to-video and video-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-InP`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-InP) | torch.float16 |
There are two official EasyAnimate checkpoints available for control-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-Control`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-Control) | torch.float16 |
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-Control-Camera`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-Control-Camera) | torch.float16 |
For the EasyAnimateV5.1 series:
- Text-to-video (T2V) and Image-to-video (I2V) works for multiple resolutions. The width and height can vary from 256 to 1024.
- Both T2V and I2V models support generation with 1~49 frames and work best at this value. Exporting videos at 8 FPS is recommended.
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`EasyAnimatePipeline`] for inference with bitsandbytes.
```py
import torch
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, EasyAnimateTransformer3DModel, EasyAnimatePipeline
from diffusers.utils import export_to_video
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True)
transformer_8bit = EasyAnimateTransformer3DModel.from_pretrained(
"alibaba-pai/EasyAnimateV5.1-12b-zh",
subfolder="transformer",
quantization_config=quant_config,
torch_dtype=torch.float16,
)
pipeline = EasyAnimatePipeline.from_pretrained(
"alibaba-pai/EasyAnimateV5.1-12b-zh",
transformer=transformer_8bit,
torch_dtype=torch.float16,
device_map="balanced",
)
prompt = "A cat walks on the grass, realistic style."
negative_prompt = "bad detailed"
video = pipeline(prompt=prompt, negative_prompt=negative_prompt, num_frames=49, num_inference_steps=30).frames[0]
export_to_video(video, "cat.mp4", fps=8)
```
## EasyAnimatePipeline
[[autodoc]] EasyAnimatePipeline
- all
- __call__
## EasyAnimatePipelineOutput
[[autodoc]] pipelines.easyanimate.pipeline_output.EasyAnimatePipelineOutput

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Flux
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Flux is a series of text-to-image generation models based on diffusion transformers. To know more about Flux, check out the original [blog post](https://blackforestlabs.ai/announcing-black-forest-labs/) by the creators of Flux, Black Forest Labs.
Original model checkpoints for Flux can be found [here](https://huggingface.co/black-forest-labs). Original inference code can be found [here](https://github.com/black-forest-labs/flux).
@@ -359,74 +355,8 @@ image.save('flux_ip_adapter_output.jpg')
<figcaption class="mt-2 text-sm text-center text-gray-500">IP-Adapter examples with prompt "wearing sunglasses"</figcaption>
</div>
## Optimize
Flux is a very large model and requires ~50GB of RAM/VRAM to load all the modeling components. Enable some of the optimizations below to lower the memory requirements.
### Group offloading
[Group offloading](../../optimization/memory#group-offloading) lowers VRAM usage by offloading groups of internal layers rather than the whole model or weights. You need to use [`~hooks.apply_group_offloading`] on all the model components of a pipeline. The `offload_type` parameter allows you to toggle between block and leaf-level offloading. Setting it to `leaf_level` offloads the lowest leaf-level parameters to the CPU instead of offloading at the module-level.
On CUDA devices that support asynchronous data streaming, set `use_stream=True` to overlap data transfer and computation to accelerate inference.
> [!TIP]
> It is possible to mix block and leaf-level offloading for different components in a pipeline.
```py
import torch
from diffusers import FluxPipeline
from diffusers.hooks import apply_group_offloading
model_id = "black-forest-labs/FLUX.1-dev"
dtype = torch.bfloat16
pipe = FluxPipeline.from_pretrained(
model_id,
torch_dtype=dtype,
)
apply_group_offloading(
pipe.transformer,
offload_type="leaf_level",
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
use_stream=True,
)
apply_group_offloading(
pipe.text_encoder,
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
offload_type="leaf_level",
use_stream=True,
)
apply_group_offloading(
pipe.text_encoder_2,
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
offload_type="leaf_level",
use_stream=True,
)
apply_group_offloading(
pipe.vae,
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
offload_type="leaf_level",
use_stream=True,
)
prompt="A cat wearing sunglasses and working as a lifeguard at pool."
generator = torch.Generator().manual_seed(181201)
image = pipe(
prompt,
width=576,
height=1024,
num_inference_steps=30,
generator=generator
).images[0]
image
```
### Running FP16 inference
## Running FP16 inference
Flux can generate high-quality images with FP16 (i.e. to accelerate inference on Turing/Volta GPUs) but produces different outputs compared to FP32/BF16. The issue is that some activations in the text encoders have to be clipped when running in FP16, which affects the overall image. Forcing text encoders to run with FP32 inference thus removes this output difference. See [here](https://github.com/huggingface/diffusers/pull/9097#issuecomment-2272292516) for details.
@@ -455,7 +385,7 @@ out = pipe(
out.save("image.png")
```
### Quantization
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.

View File

@@ -14,10 +14,6 @@
# HunyuanVideo
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[HunyuanVideo](https://www.arxiv.org/abs/2412.03603) by Tencent.
*Recent advancements in video generation have significantly impacted daily life for both individuals and industries. However, the leading video generation models remain closed-source, resulting in a notable performance gap between industry capabilities and those available to the public. In this report, we introduce HunyuanVideo, an innovative open-source video foundation model that demonstrates performance in video generation comparable to, or even surpassing, that of leading closed-source models. HunyuanVideo encompasses a comprehensive framework that integrates several key elements, including data curation, advanced architectural design, progressive model scaling and training, and an efficient infrastructure tailored for large-scale model training and inference. As a result, we successfully trained a video generative model with over 13 billion parameters, making it the largest among all open-source models. We conducted extensive experiments and implemented a series of targeted designs to ensure high visual quality, motion dynamics, text-video alignment, and advanced filming techniques. According to evaluations by professionals, HunyuanVideo outperforms previous state-of-the-art models, including Runway Gen-3, Luma 1.6, and three top-performing Chinese video generative models. By releasing the code for the foundation model and its applications, we aim to bridge the gap between closed-source and open-source communities. This initiative will empower individuals within the community to experiment with their ideas, fostering a more dynamic and vibrant video generation ecosystem. The code is publicly available at [this https URL](https://github.com/tencent/HunyuanVideo).*
@@ -36,22 +32,6 @@ Recommendations for inference:
- For smaller resolution videos, try lower values of `shift` (between `2.0` to `5.0`) in the [Scheduler](https://huggingface.co/docs/diffusers/main/en/api/schedulers/flow_match_euler_discrete#diffusers.FlowMatchEulerDiscreteScheduler.shift). For larger resolution images, try higher values (between `7.0` and `12.0`). The default value is `7.0` for HunyuanVideo.
- For more information about supported resolutions and other details, please refer to the original repository [here](https://github.com/Tencent/HunyuanVideo/).
## Available models
The following models are available for the [`HunyuanVideoPipeline`](text-to-video) pipeline:
| Model name | Description |
|:---|:---|
| [`hunyuanvideo-community/HunyuanVideo`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo) | Official HunyuanVideo (guidance-distilled). Performs best at multiple resolutions and frames. Performs best with `guidance_scale=6.0`, `true_cfg_scale=1.0` and without a negative prompt. |
| [`https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-T2V`](https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-T2V) | Skywork's custom finetune of HunyuanVideo (de-distilled). Performs best with `97x544x960` resolution, `guidance_scale=1.0`, `true_cfg_scale=6.0` and a negative prompt. |
The following models are available for the image-to-video pipeline:
| Model name | Description |
|:---|:---|
| [`Skywork/SkyReels-V1-Hunyuan-I2V`](https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-I2V) | Skywork's custom finetune of HunyuanVideo (de-distilled). Performs best with `97x544x960` resolution. Performs best at `97x544x960` resolution, `guidance_scale=1.0`, `true_cfg_scale=6.0` and a negative prompt. |
| [`hunyuanvideo-community/HunyuanVideo-I2V`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo-I2V) | Tecent's official HunyuanVideo I2V model. Performs best at resolutions of 480, 720, 960, 1280. A higher `shift` value when initializing the scheduler is recommended (good values are between 7 and 20) |
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.

View File

@@ -9,10 +9,6 @@ specific language governing permissions and limitations under the License.
# Kandinsky 3
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Kandinsky 3 is created by [Vladimir Arkhipkin](https://github.com/oriBetelgeuse),[Anastasia Maltseva](https://github.com/NastyaMittseva),[Igor Pavlov](https://github.com/boomb0om),[Andrei Filatov](https://github.com/anvilarth),[Arseniy Shakhmatov](https://github.com/cene555),[Andrey Kuznetsov](https://github.com/kuznetsoffandrey),[Denis Dimitrov](https://github.com/denndimitrov), [Zein Shaheen](https://github.com/zeinsh)
The description from it's GitHub page:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Kolors: Effective Training of Diffusion Model for Photorealistic Text-to-Image Synthesis
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/kolors/kolors_header_collage.png)
Kolors is a large-scale text-to-image generation model based on latent diffusion, developed by [the Kuaishou Kolors team](https://github.com/Kwai-Kolors/Kolors). Trained on billions of text-image pairs, Kolors exhibits significant advantages over both open-source and closed-source models in visual quality, complex semantic accuracy, and text rendering for both Chinese and English characters. Furthermore, Kolors supports both Chinese and English inputs, demonstrating strong performance in understanding and generating Chinese-specific content. For more details, please refer to this [technical report](https://github.com/Kwai-Kolors/Kolors/blob/master/imgs/Kolors_paper.pdf).

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Latent Consistency Models
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Latent Consistency Models (LCMs) were proposed in [Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference](https://huggingface.co/papers/2310.04378) by Simian Luo, Yiqin Tan, Longbo Huang, Jian Li, and Hang Zhao.
The abstract of the paper is as follows:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# LEDITS++
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
LEDITS++ was proposed in [LEDITS++: Limitless Image Editing using Text-to-Image Models](https://huggingface.co/papers/2311.16711) by Manuel Brack, Felix Friedrich, Katharina Kornmeier, Linoy Tsaban, Patrick Schramowski, Kristian Kersting, Apolinário Passos.
The abstract from the paper is:

View File

@@ -14,10 +14,6 @@
# LTX Video
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[LTX Video](https://huggingface.co/Lightricks/LTX-Video) is the first DiT-based video generation model capable of generating high-quality videos in real-time. It produces 24 FPS videos at a 768x512 resolution faster than they can be watched. Trained on a large-scale dataset of diverse videos, the model generates high-resolution videos with realistic and varied content. We provide a model for both text-to-video as well as image + text-to-video usecases.
<Tip>

View File

@@ -1,87 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# Lumina2
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Lumina Image 2.0: A Unified and Efficient Image Generative Model](https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0) is a 2 billion parameter flow-based diffusion transformer capable of generating diverse images from text descriptions.
The abstract from the paper is:
*We introduce Lumina-Image 2.0, an advanced text-to-image model that surpasses previous state-of-the-art methods across multiple benchmarks, while also shedding light on its potential to evolve into a generalist vision intelligence model. Lumina-Image 2.0 exhibits three key properties: (1) Unification it adopts a unified architecture that treats text and image tokens as a joint sequence, enabling natural cross-modal interactions and facilitating task expansion. Besides, since high-quality captioners can provide semantically better-aligned text-image training pairs, we introduce a unified captioning system, UniCaptioner, which generates comprehensive and precise captions for the model. This not only accelerates model convergence but also enhances prompt adherence, variable-length prompt handling, and task generalization via prompt templates. (2) Efficiency to improve the efficiency of the unified architecture, we develop a set of optimization techniques that improve semantic learning and fine-grained texture generation during training while incorporating inference-time acceleration strategies without compromising image quality. (3) Transparency we open-source all training details, code, and models to ensure full reproducibility, aiming to bridge the gap between well-resourced closed-source research teams and independent developers.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Using Single File loading with Lumina Image 2.0
Single file loading for Lumina Image 2.0 is available for the `Lumina2Transformer2DModel`
```python
import torch
from diffusers import Lumina2Transformer2DModel, Lumina2Text2ImgPipeline
ckpt_path = "https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0/blob/main/consolidated.00-of-01.pth"
transformer = Lumina2Transformer2DModel.from_single_file(
ckpt_path, torch_dtype=torch.bfloat16
)
pipe = Lumina2Text2ImgPipeline.from_pretrained(
"Alpha-VLLM/Lumina-Image-2.0", transformer=transformer, torch_dtype=torch.bfloat16
)
pipe.enable_model_cpu_offload()
image = pipe(
"a cat holding a sign that says hello",
generator=torch.Generator("cpu").manual_seed(0),
).images[0]
image.save("lumina-single-file.png")
```
## Using GGUF Quantized Checkpoints with Lumina Image 2.0
GGUF Quantized checkpoints for the `Lumina2Transformer2DModel` can be loaded via `from_single_file` with the `GGUFQuantizationConfig`
```python
from diffusers import Lumina2Transformer2DModel, Lumina2Text2ImgPipeline, GGUFQuantizationConfig
ckpt_path = "https://huggingface.co/calcuis/lumina-gguf/blob/main/lumina2-q4_0.gguf"
transformer = Lumina2Transformer2DModel.from_single_file(
ckpt_path,
quantization_config=GGUFQuantizationConfig(compute_dtype=torch.bfloat16),
torch_dtype=torch.bfloat16,
)
pipe = Lumina2Text2ImgPipeline.from_pretrained(
"Alpha-VLLM/Lumina-Image-2.0", transformer=transformer, torch_dtype=torch.bfloat16
)
pipe.enable_model_cpu_offload()
image = pipe(
"a cat holding a sign that says hello",
generator=torch.Generator("cpu").manual_seed(0),
).images[0]
image.save("lumina-gguf.png")
```
## Lumina2Text2ImgPipeline
[[autodoc]] Lumina2Text2ImgPipeline
- all
- __call__

View File

@@ -1,6 +1,4 @@
<!--
Copyright 2023-2025 Marigold Team, ETH Zürich. All rights reserved.
Copyright 2024-2025 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 Marigold authors and The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,120 +10,67 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Marigold Computer Vision
# Marigold Pipelines for Computer Vision Tasks
![marigold](https://marigoldmonodepth.github.io/images/teaser_collage_compressed.jpg)
Marigold was proposed in
[Repurposing Diffusion-Based Image Generators for Monocular Depth Estimation](https://huggingface.co/papers/2312.02145),
a CVPR 2024 Oral paper by
[Bingxin Ke](http://www.kebingxin.com/),
[Anton Obukhov](https://www.obukhov.ai/),
[Shengyu Huang](https://shengyuh.github.io/),
[Nando Metzger](https://nandometzger.github.io/),
[Rodrigo Caye Daudt](https://rcdaudt.github.io/), and
[Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
The core idea is to **repurpose the generative prior of Text-to-Image Latent Diffusion Models (LDMs) for traditional
computer vision tasks**.
This approach was explored by fine-tuning Stable Diffusion for **Monocular Depth Estimation**, as demonstrated in the
teaser above.
Marigold was proposed in [Repurposing Diffusion-Based Image Generators for Monocular Depth Estimation](https://huggingface.co/papers/2312.02145), a CVPR 2024 Oral paper by [Bingxin Ke](http://www.kebingxin.com/), [Anton Obukhov](https://www.obukhov.ai/), [Shengyu Huang](https://shengyuh.github.io/), [Nando Metzger](https://nandometzger.github.io/), [Rodrigo Caye Daudt](https://rcdaudt.github.io/), and [Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
The idea is to repurpose the rich generative prior of Text-to-Image Latent Diffusion Models (LDMs) for traditional computer vision tasks.
Initially, this idea was explored to fine-tune Stable Diffusion for Monocular Depth Estimation, as shown in the teaser above.
Later,
- [Tianfu Wang](https://tianfwang.github.io/) trained the first Latent Consistency Model (LCM) of Marigold, which unlocked fast single-step inference;
- [Kevin Qu](https://www.linkedin.com/in/kevin-qu-b3417621b/?locale=en_US) extended the approach to Surface Normals Estimation;
- [Anton Obukhov](https://www.obukhov.ai/) contributed the pipelines and documentation into diffusers (enabled and supported by [YiYi Xu](https://yiyixuxu.github.io/) and [Sayak Paul](https://sayak.dev/)).
Marigold was later extended in the follow-up paper,
[Marigold: Affordable Adaptation of Diffusion-Based Image Generators for Image Analysis](https://huggingface.co/papers/2312.02145),
authored by
[Bingxin Ke](http://www.kebingxin.com/),
[Kevin Qu](https://www.linkedin.com/in/kevin-qu-b3417621b/?locale=en_US),
[Tianfu Wang](https://tianfwang.github.io/),
[Nando Metzger](https://nandometzger.github.io/),
[Shengyu Huang](https://shengyuh.github.io/),
[Bo Li](https://www.linkedin.com/in/bobboli0202/),
[Anton Obukhov](https://www.obukhov.ai/), and
[Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
This work expanded Marigold to support new modalities such as **Surface Normals** and **Intrinsic Image Decomposition**
(IID), introduced a training protocol for **Latent Consistency Models** (LCM), and demonstrated **High-Resolution** (HR)
processing capability.
The abstract from the paper is:
<Tip>
The early Marigold models (`v1-0` and earlier) were optimized for best results with at least 10 inference steps.
LCM models were later developed to enable high-quality inference in just 1 to 4 steps.
Marigold models `v1-1` and later use the DDIM scheduler to achieve optimal
results in as few as 1 to 4 steps.
</Tip>
*Monocular depth estimation is a fundamental computer vision task. Recovering 3D depth from a single image is geometrically ill-posed and requires scene understanding, so it is not surprising that the rise of deep learning has led to a breakthrough. The impressive progress of monocular depth estimators has mirrored the growth in model capacity, from relatively modest CNNs to large Transformer architectures. Still, monocular depth estimators tend to struggle when presented with images with unfamiliar content and layout, since their knowledge of the visual world is restricted by the data seen during training, and challenged by zero-shot generalization to new domains. This motivates us to explore whether the extensive priors captured in recent generative diffusion models can enable better, more generalizable depth estimation. We introduce Marigold, a method for affine-invariant monocular depth estimation that is derived from Stable Diffusion and retains its rich prior knowledge. The estimator can be fine-tuned in a couple of days on a single GPU using only synthetic training data. It delivers state-of-the-art performance across a wide range of datasets, including over 20% performance gains in specific cases. Project page: https://marigoldmonodepth.github.io.*
## Available Pipelines
Each pipeline is tailored for a specific computer vision task, processing an input RGB image and generating a
corresponding prediction.
Currently, the following computer vision tasks are implemented:
Each pipeline supports one Computer Vision task, which takes an input RGB image as input and produces a *prediction* of the modality of interest, such as a depth map of the input image.
Currently, the following tasks are implemented:
| Pipeline | Predicted Modalities | Demos |
|---------------------------------------------------------------------------------------------------------------------------------------------|------------------------------------------------------------------------------------------------------------------|:--------------------------------------------------------------------------------------------------------------------------------------------------:|
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-lcm), [Slow Original Demo (DDIM)](https://huggingface.co/spaces/prs-eth/marigold) |
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-normals-lcm) |
| Pipeline | Recommended Model Checkpoints | Spaces (Interactive Apps) | Predicted Modalities |
|---------------------------------------------------------------------------------------------------------------------------------------------------|---------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|:------------------------------------------------------------------------------------:|------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [prs-eth/marigold-depth-v1-1](https://huggingface.co/prs-eth/marigold-depth-v1-1) | [Depth Estimation](https://huggingface.co/spaces/prs-eth/marigold) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) |
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [prs-eth/marigold-normals-v1-1](https://huggingface.co/prs-eth/marigold-normals-v1-1) | [Surface Normals Estimation](https://huggingface.co/spaces/prs-eth/marigold-normals) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) |
| [MarigoldIntrinsicsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_intrinsics.py) | [prs-eth/marigold-iid-appearance-v1-1](https://huggingface.co/prs-eth/marigold-iid-appearance-v1-1),<br>[prs-eth/marigold-iid-lighting-v1-1](https://huggingface.co/prs-eth/marigold-iid-lighting-v1-1) | [Intrinsic Image Decomposition](https://huggingface.co/spaces/prs-eth/marigold-iid) | [Albedo](https://en.wikipedia.org/wiki/Albedo), [Materials](https://www.n.aiq3d.com/wiki/roughnessmetalnessao-map), [Lighting](https://en.wikipedia.org/wiki/Diffuse_reflection) |
## Available Checkpoints
All original checkpoints are available under the [PRS-ETH](https://huggingface.co/prs-eth/) organization on Hugging Face.
They are designed for use with diffusers pipelines and the [original codebase](https://github.com/prs-eth/marigold), which can also be used to train
new model checkpoints.
The following is a summary of the recommended checkpoints, all of which produce reliable results with 1 to 4 steps.
| Checkpoint | Modality | Comment |
|-----------------------------------------------------------------------------------------------------|--------------|--------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
| [prs-eth/marigold-depth-v1-1](https://huggingface.co/prs-eth/marigold-depth-v1-1) | Depth | Affine-invariant depth prediction assigns each pixel a value between 0 (near plane) and 1 (far plane), with both planes determined by the model during inference. |
| [prs-eth/marigold-normals-v0-1](https://huggingface.co/prs-eth/marigold-normals-v0-1) | Normals | The surface normals predictions are unit-length 3D vectors in the screen space camera, with values in the range from -1 to 1. |
| [prs-eth/marigold-iid-appearance-v1-1](https://huggingface.co/prs-eth/marigold-iid-appearance-v1-1) | Intrinsics | InteriorVerse decomposition is comprised of Albedo and two BRDF material properties: Roughness and Metallicity. |
| [prs-eth/marigold-iid-lighting-v1-1](https://huggingface.co/prs-eth/marigold-iid-lighting-v1-1) | Intrinsics | HyperSim decomposition of an image &nbsp\\(I\\)&nbsp is comprised of Albedo &nbsp\\(A\\), Diffuse shading &nbsp\\(S\\), and Non-diffuse residual &nbsp\\(R\\): &nbsp\\(I = A*S+R\\). |
The original checkpoints can be found under the [PRS-ETH](https://huggingface.co/prs-eth/) Hugging Face organization.
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff
between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to
efficiently load the same components into multiple pipelines.
Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section
[here](../../using-diffusers/svd#reduce-memory-usage).
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines. Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section [here](../../using-diffusers/svd#reduce-memory-usage).
</Tip>
<Tip warning={true}>
Marigold pipelines were designed and tested with the scheduler embedded in the model checkpoint.
The optimal number of inference steps varies by scheduler, with no universal value that works best across all cases.
To accommodate this, the `num_inference_steps` parameter in the pipeline's `__call__` method defaults to `None` (see the
API reference).
Unless set explicitly, it inherits the value from the `default_denoising_steps` field in the checkpoint configuration
file (`model_index.json`).
This ensures high-quality predictions when invoking the pipeline with only the `image` argument.
Marigold pipelines were designed and tested only with `DDIMScheduler` and `LCMScheduler`.
Depending on the scheduler, the number of inference steps required to get reliable predictions varies, and there is no universal value that works best across schedulers.
Because of that, the default value of `num_inference_steps` in the `__call__` method of the pipeline is set to `None` (see the API reference).
Unless set explicitly, its value will be taken from the checkpoint configuration `model_index.json`.
This is done to ensure high-quality predictions when calling the pipeline with just the `image` argument.
</Tip>
See also Marigold [usage examples](../../using-diffusers/marigold_usage).
## Marigold Depth Prediction API
See also Marigold [usage examples](marigold_usage).
## MarigoldDepthPipeline
[[autodoc]] MarigoldDepthPipeline
- all
- __call__
## MarigoldNormalsPipeline
[[autodoc]] MarigoldNormalsPipeline
- all
- __call__
## MarigoldDepthOutput
[[autodoc]] pipelines.marigold.pipeline_marigold_depth.MarigoldDepthOutput
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_depth
## Marigold Normals Estimation API
[[autodoc]] MarigoldNormalsPipeline
- __call__
[[autodoc]] pipelines.marigold.pipeline_marigold_normals.MarigoldNormalsOutput
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_normals
## Marigold Intrinsic Image Decomposition API
[[autodoc]] MarigoldIntrinsicsPipeline
- __call__
[[autodoc]] pipelines.marigold.pipeline_marigold_intrinsics.MarigoldIntrinsicsOutput
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_intrinsics
## MarigoldNormalsOutput
[[autodoc]] pipelines.marigold.pipeline_marigold_normals.MarigoldNormalsOutput

View File

@@ -15,10 +15,6 @@
# Mochi 1 Preview
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
> [!TIP]
> Only a research preview of the model weights is available at the moment.

View File

@@ -1,80 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# OmniGen
[OmniGen: Unified Image Generation](https://arxiv.org/pdf/2409.11340) from BAAI, by Shitao Xiao, Yueze Wang, Junjie Zhou, Huaying Yuan, Xingrun Xing, Ruiran Yan, Chaofan Li, Shuting Wang, Tiejun Huang, Zheng Liu.
The abstract from the paper is:
*The emergence of Large Language Models (LLMs) has unified language generation tasks and revolutionized human-machine interaction. However, in the realm of image generation, a unified model capable of handling various tasks within a single framework remains largely unexplored. In this work, we introduce OmniGen, a new diffusion model for unified image generation. OmniGen is characterized by the following features: 1) Unification: OmniGen not only demonstrates text-to-image generation capabilities but also inherently supports various downstream tasks, such as image editing, subject-driven generation, and visual conditional generation. 2) Simplicity: The architecture of OmniGen is highly simplified, eliminating the need for additional plugins. Moreover, compared to existing diffusion models, it is more user-friendly and can complete complex tasks end-to-end through instructions without the need for extra intermediate steps, greatly simplifying the image generation workflow. 3) Knowledge Transfer: Benefit from learning in a unified format, OmniGen effectively transfers knowledge across different tasks, manages unseen tasks and domains, and exhibits novel capabilities. We also explore the models reasoning capabilities and potential applications of the chain-of-thought mechanism. This work represents the first attempt at a general-purpose image generation model, and we will release our resources at https://github.com/VectorSpaceLab/OmniGen to foster future advancements.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [staoxiao](https://github.com/staoxiao). The original codebase can be found [here](https://github.com/VectorSpaceLab/OmniGen). The original weights can be found under [hf.co/shitao](https://huggingface.co/Shitao/OmniGen-v1).
## Inference
First, load the pipeline:
```python
import torch
from diffusers import OmniGenPipeline
pipe = OmniGenPipeline.from_pretrained("Shitao/OmniGen-v1-diffusers", torch_dtype=torch.bfloat16)
pipe.to("cuda")
```
For text-to-image, pass a text prompt. By default, OmniGen generates a 1024x1024 image.
You can try setting the `height` and `width` parameters to generate images with different size.
```python
prompt = "Realistic photo. A young woman sits on a sofa, holding a book and facing the camera. She wears delicate silver hoop earrings adorned with tiny, sparkling diamonds that catch the light, with her long chestnut hair cascading over her shoulders. Her eyes are focused and gentle, framed by long, dark lashes. She is dressed in a cozy cream sweater, which complements her warm, inviting smile. Behind her, there is a table with a cup of water in a sleek, minimalist blue mug. The background is a serene indoor setting with soft natural light filtering through a window, adorned with tasteful art and flowers, creating a cozy and peaceful ambiance. 4K, HD."
image = pipe(
prompt=prompt,
height=1024,
width=1024,
guidance_scale=3,
generator=torch.Generator(device="cpu").manual_seed(111),
).images[0]
image.save("output.png")
```
OmniGen supports multimodal inputs.
When the input includes an image, you need to add a placeholder `<img><|image_1|></img>` in the text prompt to represent the image.
It is recommended to enable `use_input_image_size_as_output` to keep the edited image the same size as the original image.
```python
prompt="<img><|image_1|></img> Remove the woman's earrings. Replace the mug with a clear glass filled with sparkling iced cola."
input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png")]
image = pipe(
prompt=prompt,
input_images=input_images,
guidance_scale=2,
img_guidance_scale=1.6,
use_input_image_size_as_output=True,
generator=torch.Generator(device="cpu").manual_seed(222)).images[0]
image.save("output.png")
```
## OmniGenPipeline
[[autodoc]] OmniGenPipeline
- all
- __call__

View File

@@ -54,7 +54,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [DiT](dit) | text2image |
| [Flux](flux) | text2image |
| [Hunyuan-DiT](hunyuandit) | text2image |
| [I2VGen-XL](i2vgenxl) | image2video |
| [I2VGen-XL](i2vgenxl) | text2video |
| [InstructPix2Pix](pix2pix) | image editing |
| [Kandinsky 2.1](kandinsky) | text2image, image2image, inpainting, interpolation |
| [Kandinsky 2.2](kandinsky_v22) | text2image, image2image, inpainting |
@@ -65,7 +65,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Latte](latte) | text2image |
| [LEDITS++](ledits_pp) | image editing |
| [Lumina-T2X](lumina) | text2image |
| [Marigold](marigold) | depth-estimation, normals-estimation, intrinsic-decomposition |
| [Marigold](marigold) | depth |
| [MultiDiffusion](panorama) | text2image |
| [MusicLDM](musicldm) | text2audio |
| [PAG](pag) | text2image |

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Perturbed-Attention Guidance
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Perturbed-Attention Guidance (PAG)](https://ku-cvlab.github.io/Perturbed-Attention-Guidance/) is a new diffusion sampling guidance that improves sample quality across both unconditional and conditional settings, achieving this without requiring further training or the integration of external modules.
PAG was introduced in [Self-Rectifying Diffusion Sampling with Perturbed-Attention Guidance](https://huggingface.co/papers/2403.17377) by Donghoon Ahn, Hyoungwon Cho, Jaewon Min, Wooseok Jang, Jungwoo Kim, SeonHwa Kim, Hyun Hee Park, Kyong Hwan Jin and Seungryong Kim.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# MultiDiffusion
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation](https://huggingface.co/papers/2302.08113) is by Omer Bar-Tal, Lior Yariv, Yaron Lipman, and Tali Dekel.
The abstract from the paper is:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Image-to-Video Generation with PIA (Personalized Image Animator)
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
## Overview
[PIA: Your Personalized Image Animator via Plug-and-Play Modules in Text-to-Image Models](https://arxiv.org/abs/2312.13964) by Yiming Zhang, Zhening Xing, Yanhong Zeng, Youqing Fang, Kai Chen

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# InstructPix2Pix
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[InstructPix2Pix: Learning to Follow Image Editing Instructions](https://huggingface.co/papers/2211.09800) is by Tim Brooks, Aleksander Holynski and Alexei A. Efros.
The abstract from the paper is:

View File

@@ -14,10 +14,6 @@
# SanaPipeline
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[SANA: Efficient High-Resolution Image Synthesis with Linear Diffusion Transformers](https://huggingface.co/papers/2410.10629) from NVIDIA and MIT HAN Lab, by Enze Xie, Junsong Chen, Junyu Chen, Han Cai, Haotian Tang, Yujun Lin, Zhekai Zhang, Muyang Li, Ligeng Zhu, Yao Lu, Song Han.
The abstract from the paper is:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Depth-to-image
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model can also infer depth based on an image using [MiDaS](https://github.com/isl-org/MiDaS). This allows you to pass a text prompt and an initial image to condition the generation of new images as well as a `depth_map` to preserve the image structure.
<Tip>

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Image-to-image
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model can also be applied to image-to-image generation by passing a text prompt and an initial image to condition the generation of new images.
The [`StableDiffusionImg2ImgPipeline`] uses the diffusion-denoising mechanism proposed in [SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equations](https://huggingface.co/papers/2108.01073) by Chenlin Meng, Yutong He, Yang Song, Jiaming Song, Jiajun Wu, Jun-Yan Zhu, Stefano Ermon.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Inpainting
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model can also be applied to inpainting which lets you edit specific parts of an image by providing a mask and a text prompt using Stable Diffusion.
## Tips

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text-to-(RGB, depth)
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
LDM3D was proposed in [LDM3D: Latent Diffusion Model for 3D](https://huggingface.co/papers/2305.10853) by Gabriela Ben Melech Stan, Diana Wofk, Scottie Fox, Alex Redden, Will Saxton, Jean Yu, Estelle Aflalo, Shao-Yen Tseng, Fabio Nonato, Matthias Muller, and Vasudev Lal. LDM3D generates an image and a depth map from a given text prompt unlike the existing text-to-image diffusion models such as [Stable Diffusion](./overview) which only generates an image. With almost the same number of parameters, LDM3D achieves to create a latent space that can compress both the RGB images and the depth maps.
Two checkpoints are available for use:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable Diffusion pipelines
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/). Latent diffusion applies the diffusion process over a lower dimensional latent space to reduce memory and compute complexity. This specific type of diffusion model was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
Stable Diffusion is trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable Diffusion 3
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Stable Diffusion 3 (SD3) was proposed in [Scaling Rectified Flow Transformers for High-Resolution Image Synthesis](https://arxiv.org/pdf/2403.03206.pdf) by Patrick Esser, Sumith Kulal, Andreas Blattmann, Rahim Entezari, Jonas Muller, Harry Saini, Yam Levi, Dominik Lorenz, Axel Sauer, Frederic Boesel, Dustin Podell, Tim Dockhorn, Zion English, Kyle Lacey, Alex Goodwin, Yannik Marek, and Robin Rombach.
The abstract from the paper is:
@@ -81,7 +77,7 @@ from diffusers import StableDiffusion3Pipeline
from transformers import SiglipVisionModel, SiglipImageProcessor
image_encoder_id = "google/siglip-so400m-patch14-384"
ip_adapter_id = "InstantX/SD3.5-Large-IP-Adapter"
ip_adapter_id = "guiyrt/InstantX-SD3.5-Large-IP-Adapter-diffusers"
feature_extractor = SiglipImageProcessor.from_pretrained(
image_encoder_id,

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable Diffusion XL
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Stable Diffusion XL (SDXL) was proposed in [SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis](https://huggingface.co/papers/2307.01952) by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, and Robin Rombach.
The abstract from the paper is:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text-to-image
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model was created by researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [Runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionPipeline`] is capable of generating photorealistic images given any text input. It's trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs. Latent diffusion is the research on top of which Stable Diffusion was built. It was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
The abstract from the paper is:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Super-resolution
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion upscaler diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), and [LAION](https://laion.ai/). It is used to enhance the resolution of input images by a factor of 4.
<Tip>

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable unCLIP
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Stable unCLIP checkpoints are finetuned from [Stable Diffusion 2.1](./stable_diffusion/stable_diffusion_2) checkpoints to condition on CLIP image embeddings.
Stable unCLIP still conditions on text embeddings. Given the two separate conditionings, stable unCLIP can be used
for text guided image variation. When combined with an unCLIP prior, it can also be used for full text to image generation.

View File

@@ -18,10 +18,6 @@ specific language governing permissions and limitations under the License.
# Text-to-video
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[ModelScope Text-to-Video Technical Report](https://arxiv.org/abs/2308.06571) is by Jiuniu Wang, Hangjie Yuan, Dayou Chen, Yingya Zhang, Xiang Wang, Shiwei Zhang.
The abstract from the paper is:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text2Video-Zero
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Text2Video-Zero: Text-to-Image Diffusion Models are Zero-Shot Video Generators](https://huggingface.co/papers/2303.13439) is by Levon Khachatryan, Andranik Movsisyan, Vahram Tadevosyan, Roberto Henschel, [Zhangyang Wang](https://www.ece.utexas.edu/people/faculty/atlas-wang), Shant Navasardyan, [Humphrey Shi](https://www.humphreyshi.com).
Text2Video-Zero enables zero-shot video generation using either:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# UniDiffuser
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The UniDiffuser model was proposed in [One Transformer Fits All Distributions in Multi-Modal Diffusion at Scale](https://huggingface.co/papers/2303.06555) by Fan Bao, Shen Nie, Kaiwen Xue, Chongxuan Li, Shi Pu, Yaole Wang, Gang Yue, Yue Cao, Hang Su, Jun Zhu.
The abstract from the paper is:

View File

@@ -1,78 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# Wan
[Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
<!-- TODO(aryan): update abstract once paper is out -->
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
Recommendations for inference:
- VAE in `torch.float32` for better decoding quality.
- `num_frames` should be of the form `4 * k + 1`, for example `49` or `81`.
- For smaller resolution videos, try lower values of `shift` (between `2.0` to `5.0`) in the [Scheduler](https://huggingface.co/docs/diffusers/main/en/api/schedulers/flow_match_euler_discrete#diffusers.FlowMatchEulerDiscreteScheduler.shift). For larger resolution videos, try higher values (between `7.0` and `12.0`). The default value is `3.0` for Wan.
### Using a custom scheduler
Wan can be used with many different schedulers, each with their own benefits regarding speed and generation quality. By default, Wan uses the `UniPCMultistepScheduler(prediction_type="flow_prediction", use_flow_sigmas=True, flow_shift=3.0)` scheduler. You can use a different scheduler as follows:
```python
from diffusers import FlowMatchEulerDiscreteScheduler, UniPCMultistepScheduler, WanPipeline
scheduler_a = FlowMatchEulerDiscreteScheduler(shift=5.0)
scheduler_b = UniPCMultistepScheduler(prediction_type="flow_prediction", use_flow_sigmas=True, flow_shift=4.0)
pipe = WanPipeline.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", scheduler=<CUSTOM_SCHEDULER_HERE>)
# or,
pipe.scheduler = <CUSTOM_SCHEDULER_HERE>
```
### Using single file loading with Wan
The `WanTransformer3DModel` and `AutoencoderKLWan` models support loading checkpoints in their original format via the `from_single_file` loading
method.
```python
import torch
from diffusers import WanPipeline, WanTransformer3DModel
ckpt_path = "https://huggingface.co/Comfy-Org/Wan_2.1_ComfyUI_repackaged/blob/main/split_files/diffusion_models/wan2.1_t2v_1.3B_bf16.safetensors"
transformer = WanTransformer3DModel.from_single_file(ckpt_path, torch_dtype=torch.bfloat16)
pipe = WanPipeline.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", transformer=transformer)
```
## WanPipeline
[[autodoc]] WanPipeline
- all
- __call__
## WanImageToVideoPipeline
[[autodoc]] WanImageToVideoPipeline
- all
- __call__
## WanPipelineOutput
[[autodoc]] pipelines.wan.pipeline_output.WanPipelineOutput

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Würstchen
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
<img src="https://github.com/dome272/Wuerstchen/assets/61938694/0617c863-165a-43ee-9303-2a17299a0cf9">
[Wuerstchen: An Efficient Architecture for Large-Scale Text-to-Image Diffusion Models](https://huggingface.co/papers/2306.00637) is by Pablo Pernias, Dominic Rampas, Mats L. Richter and Christopher Pal and Marc Aubreville.

View File

@@ -31,11 +31,6 @@ Learn how to quantize models in the [Quantization](../quantization/overview) gui
## GGUFQuantizationConfig
[[autodoc]] GGUFQuantizationConfig
## QuantoConfig
[[autodoc]] QuantoConfig
## TorchAoConfig
[[autodoc]] TorchAoConfig

View File

@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# CogVideoXDDIMScheduler
`CogVideoXDDIMScheduler` is based on [Denoising Diffusion Implicit Models](https://huggingface.co/papers/2010.02502), specifically for CogVideoX models.
## CogVideoXDDIMScheduler
[[autodoc]] CogVideoXDDIMScheduler

View File

@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# CogVideoXDPMScheduler
`CogVideoXDPMScheduler` is based on [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095), specifically for CogVideoX models.
## CogVideoXDPMScheduler
[[autodoc]] CogVideoXDPMScheduler

View File

@@ -45,7 +45,3 @@ Utility and helper functions for working with 🤗 Diffusers.
## apply_layerwise_casting
[[autodoc]] hooks.layerwise_casting.apply_layerwise_casting
## apply_group_offloading
[[autodoc]] hooks.group_offloading.apply_group_offloading

View File

@@ -16,11 +16,6 @@ specific language governing permissions and limitations under the License.
<img src="https://colab.research.google.com/assets/colab-badge.svg" alt="Open In Colab"/>
</a>
> [!TIP]
> This document has now grown outdated given the emergence of existing evaluation frameworks for diffusion models for image generation. Please check
> out works like [HEIM](https://crfm.stanford.edu/helm/heim/latest/), [T2I-Compbench](https://arxiv.org/abs/2307.06350),
> [GenEval](https://arxiv.org/abs/2310.11513).
Evaluation of generative models like [Stable Diffusion](https://huggingface.co/docs/diffusers/stable_diffusion) is subjective in nature. But as practitioners and researchers, we often have to make careful choices amongst many different possibilities. So, when working with different generative models (like GANs, Diffusion, etc.), how do we choose one over the other?
Qualitative evaluation of such models can be error-prone and might incorrectly influence a decision.

View File

@@ -1,5 +0,0 @@
# Hybrid Inference API Reference
## Remote Decode
[[autodoc]] utils.remote_utils.remote_decode

View File

@@ -1,54 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Hybrid Inference
**Empowering local AI builders with Hybrid Inference**
> [!TIP]
> Hybrid Inference is an [experimental feature](https://huggingface.co/blog/remote_vae).
> Feedback can be provided [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
## Why use Hybrid Inference?
Hybrid Inference offers a fast and simple way to offload local generation requirements.
- 🚀 **Reduced Requirements:** Access powerful models without expensive hardware.
- 💎 **Without Compromise:** Achieve the highest quality without sacrificing performance.
- 💰 **Cost Effective:** It's free! 🤑
- 🎯 **Diverse Use Cases:** Fully compatible with Diffusers 🧨 and the wider community.
- 🔧 **Developer-Friendly:** Simple requests, fast responses.
---
## Available Models
* **VAE Decode 🖼️:** Quickly decode latent representations into high-quality images without compromising performance or workflow speed.
* **VAE Encode 🔢 (coming soon):** Efficiently encode images into latent representations for generation and training.
* **Text Encoders 📃 (coming soon):** Compute text embeddings for your prompts quickly and accurately, ensuring a smooth and high-quality workflow.
---
## Integrations
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.
## Contents
The documentation is organized into two sections:
* **VAE Decode** Learn the basics of how to use VAE Decode with Hybrid Inference.
* **API Reference** Dive into task-specific settings and parameters.

View File

@@ -1,345 +0,0 @@
# Getting Started: VAE Decode with Hybrid Inference
VAE decode is an essential component of diffusion models - turning latent representations into images or videos.
## Memory
These tables demonstrate the VRAM requirements for VAE decode with SD v1 and SD XL on different GPUs.
For the majority of these GPUs the memory usage % dictates other models (text encoders, UNet/Transformer) must be offloaded, or tiled decoding has to be used which increases time taken and impacts quality.
<details><summary>SD v1.5</summary>
| GPU | Resolution | Time (seconds) | Memory (%) | Tiled Time (secs) | Tiled Memory (%) |
| --- | --- | --- | --- | --- | --- |
| NVIDIA GeForce RTX 4090 | 512x512 | 0.031 | 5.60% | 0.031 (0%) | 5.60% |
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.148 | 20.00% | 0.301 (+103%) | 5.60% |
| NVIDIA GeForce RTX 4080 | 512x512 | 0.05 | 8.40% | 0.050 (0%) | 8.40% |
| NVIDIA GeForce RTX 4080 | 1024x1024 | 0.224 | 30.00% | 0.356 (+59%) | 8.40% |
| NVIDIA GeForce RTX 4070 Ti | 512x512 | 0.066 | 11.30% | 0.066 (0%) | 11.30% |
| NVIDIA GeForce RTX 4070 Ti | 1024x1024 | 0.284 | 40.50% | 0.454 (+60%) | 11.40% |
| NVIDIA GeForce RTX 3090 | 512x512 | 0.062 | 5.20% | 0.062 (0%) | 5.20% |
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.253 | 18.50% | 0.464 (+83%) | 5.20% |
| NVIDIA GeForce RTX 3080 | 512x512 | 0.07 | 12.80% | 0.070 (0%) | 12.80% |
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.286 | 45.30% | 0.466 (+63%) | 12.90% |
| NVIDIA GeForce RTX 3070 | 512x512 | 0.102 | 15.90% | 0.102 (0%) | 15.90% |
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.421 | 56.30% | 0.746 (+77%) | 16.00% |
</details>
<details><summary>SDXL</summary>
| GPU | Resolution | Time (seconds) | Memory Consumed (%) | Tiled Time (seconds) | Tiled Memory (%) |
| --- | --- | --- | --- | --- | --- |
| NVIDIA GeForce RTX 4090 | 512x512 | 0.057 | 10.00% | 0.057 (0%) | 10.00% |
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.256 | 35.50% | 0.257 (+0.4%) | 35.50% |
| NVIDIA GeForce RTX 4080 | 512x512 | 0.092 | 15.00% | 0.092 (0%) | 15.00% |
| NVIDIA GeForce RTX 4080 | 1024x1024 | 0.406 | 53.30% | 0.406 (0%) | 53.30% |
| NVIDIA GeForce RTX 4070 Ti | 512x512 | 0.121 | 20.20% | 0.120 (-0.8%) | 20.20% |
| NVIDIA GeForce RTX 4070 Ti | 1024x1024 | 0.519 | 72.00% | 0.519 (0%) | 72.00% |
| NVIDIA GeForce RTX 3090 | 512x512 | 0.107 | 10.50% | 0.107 (0%) | 10.50% |
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.459 | 38.00% | 0.460 (+0.2%) | 38.00% |
| NVIDIA GeForce RTX 3080 | 512x512 | 0.121 | 25.60% | 0.121 (0%) | 25.60% |
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.524 | 93.00% | 0.524 (0%) | 93.00% |
| NVIDIA GeForce RTX 3070 | 512x512 | 0.183 | 31.80% | 0.183 (0%) | 31.80% |
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.794 | 96.40% | 0.794 (0%) | 96.40% |
</details>
## Available VAEs
| | **Endpoint** | **Model** |
|:-:|:-----------:|:--------:|
| **Stable Diffusion v1** | [https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud](https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud) | [`stabilityai/sd-vae-ft-mse`](https://hf.co/stabilityai/sd-vae-ft-mse) |
| **Stable Diffusion XL** | [https://x2dmsqunjd6k9prw.us-east-1.aws.endpoints.huggingface.cloud](https://x2dmsqunjd6k9prw.us-east-1.aws.endpoints.huggingface.cloud) | [`madebyollin/sdxl-vae-fp16-fix`](https://hf.co/madebyollin/sdxl-vae-fp16-fix) |
| **Flux** | [https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud](https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud) | [`black-forest-labs/FLUX.1-schnell`](https://hf.co/black-forest-labs/FLUX.1-schnell) |
| **HunyuanVideo** | [https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud](https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud) | [`hunyuanvideo-community/HunyuanVideo`](https://hf.co/hunyuanvideo-community/HunyuanVideo) |
> [!TIP]
> Model support can be requested [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
## Code
> [!TIP]
> Install `diffusers` from `main` to run the code: `pip install git+https://github.com/huggingface/diffusers@main`
A helper method simplifies interacting with Hybrid Inference.
```python
from diffusers.utils.remote_utils import remote_decode
```
### Basic example
Here, we show how to use the remote VAE on random tensors.
<details><summary>Code</summary>
```python
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=torch.randn([1, 4, 64, 64], dtype=torch.float16),
scaling_factor=0.18215,
)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/output.png"/>
</figure>
Usage for Flux is slightly different. Flux latents are packed so we need to send the `height` and `width`.
<details><summary>Code</summary>
```python
image = remote_decode(
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=torch.randn([1, 4096, 64], dtype=torch.float16),
height=1024,
width=1024,
scaling_factor=0.3611,
shift_factor=0.1159,
)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/flux_random_latent.png"/>
</figure>
Finally, an example for HunyuanVideo.
<details><summary>Code</summary>
```python
video = remote_decode(
endpoint="https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=torch.randn([1, 16, 3, 40, 64], dtype=torch.float16),
output_type="mp4",
)
with open("video.mp4", "wb") as f:
f.write(video)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<video
alt="queue.mp4"
autoplay loop autobuffer muted playsinline
>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/video_1.mp4" type="video/mp4">
</video>
</figure>
### Generation
But we want to use the VAE on an actual pipeline to get an actual image, not random noise. The example below shows how to do it with SD v1.5.
<details><summary>Code</summary>
```python
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
variant="fp16",
vae=None,
).to("cuda")
prompt = "Strawberry ice cream, in a stylish modern glass, coconut, splashing milk cream and honey, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious"
latent = pipe(
prompt=prompt,
output_type="latent",
).images
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
scaling_factor=0.18215,
)
image.save("test.jpg")
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/test.jpg"/>
</figure>
Heres another example with Flux.
<details><summary>Code</summary>
```python
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained(
"black-forest-labs/FLUX.1-schnell",
torch_dtype=torch.bfloat16,
vae=None,
).to("cuda")
prompt = "Strawberry ice cream, in a stylish modern glass, coconut, splashing milk cream and honey, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious"
latent = pipe(
prompt=prompt,
guidance_scale=0.0,
num_inference_steps=4,
output_type="latent",
).images
image = remote_decode(
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
height=1024,
width=1024,
scaling_factor=0.3611,
shift_factor=0.1159,
)
image.save("test.jpg")
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/test_1.jpg"/>
</figure>
Heres an example with HunyuanVideo.
<details><summary>Code</summary>
```python
from diffusers import HunyuanVideoPipeline, HunyuanVideoTransformer3DModel
model_id = "hunyuanvideo-community/HunyuanVideo"
transformer = HunyuanVideoTransformer3DModel.from_pretrained(
model_id, subfolder="transformer", torch_dtype=torch.bfloat16
)
pipe = HunyuanVideoPipeline.from_pretrained(
model_id, transformer=transformer, vae=None, torch_dtype=torch.float16
).to("cuda")
latent = pipe(
prompt="A cat walks on the grass, realistic",
height=320,
width=512,
num_frames=61,
num_inference_steps=30,
output_type="latent",
).frames
video = remote_decode(
endpoint="https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
output_type="mp4",
)
if isinstance(video, bytes):
with open("video.mp4", "wb") as f:
f.write(video)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<video
alt="queue.mp4"
autoplay loop autobuffer muted playsinline
>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/video.mp4" type="video/mp4">
</video>
</figure>
### Queueing
One of the great benefits of using a remote VAE is that we can queue multiple generation requests. While the current latent is being processed for decoding, we can already queue another one. This helps improve concurrency.
<details><summary>Code</summary>
```python
import queue
import threading
from IPython.display import display
from diffusers import StableDiffusionPipeline
def decode_worker(q: queue.Queue):
while True:
item = q.get()
if item is None:
break
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=item,
scaling_factor=0.18215,
)
display(image)
q.task_done()
q = queue.Queue()
thread = threading.Thread(target=decode_worker, args=(q,), daemon=True)
thread.start()
def decode(latent: torch.Tensor):
q.put(latent)
prompts = [
"Blueberry ice cream, in a stylish modern glass , ice cubes, nuts, mint leaves, splashing milk cream, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious",
"Lemonade in a glass, mint leaves, in an aqua and white background, flowers, ice cubes, halo, fluid motion, dynamic movement, soft lighting, digital painting, rule of thirds composition, Art by Greg rutkowski, Coby whitmore",
"Comic book art, beautiful, vintage, pastel neon colors, extremely detailed pupils, delicate features, light on face, slight smile, Artgerm, Mary Blair, Edmund Dulac, long dark locks, bangs, glowing, fashionable style, fairytale ambience, hot pink.",
"Masterpiece, vanilla cone ice cream garnished with chocolate syrup, crushed nuts, choco flakes, in a brown background, gold, cinematic lighting, Art by WLOP",
"A bowl of milk, falling cornflakes, berries, blueberries, in a white background, soft lighting, intricate details, rule of thirds, octane render, volumetric lighting",
"Cold Coffee with cream, crushed almonds, in a glass, choco flakes, ice cubes, wet, in a wooden background, cinematic lighting, hyper realistic painting, art by Carne Griffiths, octane render, volumetric lighting, fluid motion, dynamic movement, muted colors,",
]
pipe = StableDiffusionPipeline.from_pretrained(
"Lykon/dreamshaper-8",
torch_dtype=torch.float16,
vae=None,
).to("cuda")
pipe.unet = pipe.unet.to(memory_format=torch.channels_last)
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
_ = pipe(
prompt=prompts[0],
output_type="latent",
)
for prompt in prompts:
latent = pipe(
prompt=prompt,
output_type="latent",
).images
decode(latent)
q.put(None)
thread.join()
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<video
alt="queue.mp4"
autoplay loop autobuffer muted playsinline
>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/queue.mp4" type="video/mp4">
</video>
</figure>
## Integrations
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.

View File

@@ -158,46 +158,6 @@ In order to properly offload models after they're called, it is required to run
</Tip>
## Group offloading
Group offloading is the middle ground between sequential and model offloading. It works by offloading groups of internal layers (either `torch.nn.ModuleList` or `torch.nn.Sequential`), which uses less memory than model-level offloading. It is also faster than sequential-level offloading because the number of device synchronizations is reduced.
To enable group offloading, call the [`~ModelMixin.enable_group_offload`] method on the model if it is a Diffusers model implementation. For any other model implementation, use [`~hooks.group_offloading.apply_group_offloading`]:
```python
import torch
from diffusers import CogVideoXPipeline
from diffusers.hooks import apply_group_offloading
from diffusers.utils import export_to_video
# Load the pipeline
onload_device = torch.device("cuda")
offload_device = torch.device("cpu")
pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
# We can utilize the enable_group_offload method for Diffusers model implementations
pipe.transformer.enable_group_offload(onload_device=onload_device, offload_device=offload_device, offload_type="leaf_level", use_stream=True)
# For any other model implementations, the apply_group_offloading function can be used
apply_group_offloading(pipe.text_encoder, onload_device=onload_device, offload_type="block_level", num_blocks_per_group=2)
apply_group_offloading(pipe.vae, onload_device=onload_device, offload_type="leaf_level")
prompt = (
"A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. "
"The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other "
"pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, "
"casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. "
"The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical "
"atmosphere of this unique musical performance."
)
video = pipe(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
# This utilized about 14.79 GB. It can be further reduced by using tiling and using leaf_level offloading throughout the pipeline.
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
export_to_video(video, "output.mp4", fps=8)
```
Group offloading (for CUDA devices with support for asynchronous data transfer streams) overlaps data transfer and computation to reduce the overall execution time compared to sequential offloading. This is enabled using layer prefetching with CUDA streams. The next layer to be executed is loaded onto the accelerator device while the current layer is being executed - this increases the memory requirements slightly. Group offloading also supports leaf-level offloading (equivalent to sequential CPU offloading) but can be made much faster when using streams.
## FP8 layerwise weight-casting
PyTorch supports `torch.float8_e4m3fn` and `torch.float8_e5m2` as weight storage dtypes, but they can't be used for computation in many different tensor operations due to unimplemented kernel support. However, you can use these dtypes to store model weights in fp8 precision and upcast them on-the-fly when the layers are used in the forward pass. This is known as layerwise weight-casting.

View File

@@ -36,6 +36,5 @@ Diffusers currently supports the following quantization methods.
- [BitsandBytes](./bitsandbytes)
- [TorchAO](./torchao)
- [GGUF](./gguf)
- [Quanto](./quanto.md)
[This resource](https://huggingface.co/docs/transformers/main/en/quantization/overview#when-to-use-what) provides a good overview of the pros and cons of different quantization techniques.

View File

@@ -1,148 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Quanto
[Quanto](https://github.com/huggingface/optimum-quanto) is a PyTorch quantization backend for [Optimum](https://huggingface.co/docs/optimum/en/index). It has been designed with versatility and simplicity in mind:
- All features are available in eager mode (works with non-traceable models)
- Supports quantization aware training
- Quantized models are compatible with `torch.compile`
- Quantized models are Device agnostic (e.g CUDA,XPU,MPS,CPU)
In order to use the Quanto backend, you will first need to install `optimum-quanto>=0.2.6` and `accelerate`
```shell
pip install optimum-quanto accelerate
```
Now you can quantize a model by passing the `QuantoConfig` object to the `from_pretrained()` method. Although the Quanto library does allow quantizing `nn.Conv2d` and `nn.LayerNorm` modules, currently, Diffusers only supports quantizing the weights in the `nn.Linear` layers of a model. The following snippet demonstrates how to apply `float8` quantization with Quanto.
```python
import torch
from diffusers import FluxTransformer2DModel, QuantoConfig
model_id = "black-forest-labs/FLUX.1-dev"
quantization_config = QuantoConfig(weights_dtype="float8")
transformer = FluxTransformer2DModel.from_pretrained(
model_id,
subfolder="transformer",
quantization_config=quantization_config,
torch_dtype=torch.bfloat16,
)
pipe = FluxPipeline.from_pretrained(model_id, transformer=transformer, torch_dtype=torch_dtype)
pipe.to("cuda")
prompt = "A cat holding a sign that says hello world"
image = pipe(
prompt, num_inference_steps=50, guidance_scale=4.5, max_sequence_length=512
).images[0]
image.save("output.png")
```
## Skipping Quantization on specific modules
It is possible to skip applying quantization on certain modules using the `modules_to_not_convert` argument in the `QuantoConfig`. Please ensure that the modules passed in to this argument match the keys of the modules in the `state_dict`
```python
import torch
from diffusers import FluxTransformer2DModel, QuantoConfig
model_id = "black-forest-labs/FLUX.1-dev"
quantization_config = QuantoConfig(weights_dtype="float8", modules_to_not_convert=["proj_out"])
transformer = FluxTransformer2DModel.from_pretrained(
model_id,
subfolder="transformer",
quantization_config=quantization_config,
torch_dtype=torch.bfloat16,
)
```
## Using `from_single_file` with the Quanto Backend
`QuantoConfig` is compatible with `~FromOriginalModelMixin.from_single_file`.
```python
import torch
from diffusers import FluxTransformer2DModel, QuantoConfig
ckpt_path = "https://huggingface.co/black-forest-labs/FLUX.1-dev/blob/main/flux1-dev.safetensors"
quantization_config = QuantoConfig(weights_dtype="float8")
transformer = FluxTransformer2DModel.from_single_file(ckpt_path, quantization_config=quantization_config, torch_dtype=torch.bfloat16)
```
## Saving Quantized models
Diffusers supports serializing Quanto models using the `~ModelMixin.save_pretrained` method.
The serialization and loading requirements are different for models quantized directly with the Quanto library and models quantized
with Diffusers using Quanto as the backend. It is currently not possible to load models quantized directly with Quanto into Diffusers using `~ModelMixin.from_pretrained`
```python
import torch
from diffusers import FluxTransformer2DModel, QuantoConfig
model_id = "black-forest-labs/FLUX.1-dev"
quantization_config = QuantoConfig(weights_dtype="float8")
transformer = FluxTransformer2DModel.from_pretrained(
model_id,
subfolder="transformer",
quantization_config=quantization_config,
torch_dtype=torch.bfloat16,
)
# save quantized model to reuse
transformer.save_pretrained("<your quantized model save path>")
# you can reload your quantized model with
model = FluxTransformer2DModel.from_pretrained("<your quantized model save path>")
```
## Using `torch.compile` with Quanto
Currently the Quanto backend supports `torch.compile` for the following quantization types:
- `int8` weights
```python
import torch
from diffusers import FluxPipeline, FluxTransformer2DModel, QuantoConfig
model_id = "black-forest-labs/FLUX.1-dev"
quantization_config = QuantoConfig(weights_dtype="int8")
transformer = FluxTransformer2DModel.from_pretrained(
model_id,
subfolder="transformer",
quantization_config=quantization_config,
torch_dtype=torch.bfloat16,
)
transformer = torch.compile(transformer, mode="max-autotune", fullgraph=True)
pipe = FluxPipeline.from_pretrained(
model_id, transformer=transformer, torch_dtype=torch_dtype
)
pipe.to("cuda")
images = pipe("A cat holding a sign that says hello").images[0]
images.save("flux-quanto-compile.png")
```
## Supported Quantization Types
### Weights
- float8
- int8
- int4
- int2

View File

@@ -126,7 +126,7 @@ image = pipe(prompt, num_inference_steps=30, guidance_scale=7.0).images[0]
image.save("output.png")
```
If you are using `torch<=2.6.0`, some quantization methods, such as `uint4wo`, cannot be loaded directly and may result in an `UnpicklingError` when trying to load the models, but work as expected when saving them. In order to work around this, one can load the state dict manually into the model. Note, however, that this requires using `weights_only=False` in `torch.load`, so it should be run only if the weights were obtained from a trustable source.
Some quantization methods, such as `uint4wo`, cannot be loaded directly and may result in an `UnpicklingError` when trying to load the models, but work as expected when saving them. In order to work around this, one can load the state dict manually into the model. Note, however, that this requires using `weights_only=False` in `torch.load`, so it should be run only if the weights were obtained from a trustable source.
```python
import torch

View File

@@ -339,10 +339,7 @@ import torch
from huggingface_hub.repocard import RepoCard
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16,
).to("cuda")
model_id = "sayakpaul/custom-diffusion-cat-wooden-pot"
pipeline = DiffusionPipeline.from_pretrained("sayakpaul/custom-diffusion-cat-wooden-pot", torch_dtype=torch.float16).to("cuda")
pipeline.unet.load_attn_procs(model_id, weight_name="pytorch_custom_diffusion_weights.bin")
pipeline.load_textual_inversion(model_id, weight_name="<new1>.bin")
pipeline.load_textual_inversion(model_id, weight_name="<new2>.bin")

View File

@@ -221,7 +221,3 @@ pipe.delete_adapters("toy")
pipe.get_active_adapters()
["pixel"]
```
## PeftInputAutocastDisableHook
[[autodoc]] hooks.layerwise_casting.PeftInputAutocastDisableHook

View File

@@ -157,84 +157,6 @@ pipeline(
)
```
## IP Adapter Cutoff
IP Adapter is an image prompt adapter that can be used for diffusion models without any changes to the underlying model. We can use the IP Adapter Cutoff Callback to disable the IP Adapter after a certain number of steps. To set up the callback, you need to specify the number of denoising steps after which the callback comes into effect. You can do so by using either one of these two arguments:
- `cutoff_step_ratio`: Float number with the ratio of the steps.
- `cutoff_step_index`: Integer number with the exact number of the step.
We need to download the diffusion model and load the ip_adapter for it as follows:
```py
from diffusers import AutoPipelineForText2Image
from diffusers.utils import load_image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16).to("cuda")
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="sdxl_models", weight_name="ip-adapter_sdxl.bin")
pipeline.set_ip_adapter_scale(0.6)
```
The setup for the callback should look something like this:
```py
from diffusers import AutoPipelineForText2Image
from diffusers.callbacks import IPAdapterScaleCutoffCallback
from diffusers.utils import load_image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16
).to("cuda")
pipeline.load_ip_adapter(
"h94/IP-Adapter",
subfolder="sdxl_models",
weight_name="ip-adapter_sdxl.bin"
)
pipeline.set_ip_adapter_scale(0.6)
callback = IPAdapterScaleCutoffCallback(
cutoff_step_ratio=None,
cutoff_step_index=5
)
image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_adapter_diner.png"
)
generator = torch.Generator(device="cuda").manual_seed(2628670641)
images = pipeline(
prompt="a tiger sitting in a chair drinking orange juice",
ip_adapter_image=image,
negative_prompt="deformed, ugly, wrong proportion, low res, bad anatomy, worst quality, low quality",
generator=generator,
num_inference_steps=50,
callback_on_step_end=callback,
).images
images[0].save("custom_callback_img.png")
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/without_callback.png" alt="generated image of a tiger sitting in a chair drinking orange juice" />
<figcaption class="mt-2 text-center text-sm text-gray-500">without IPAdapterScaleCutoffCallback</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/with_callback2.png" alt="generated image of a tiger sitting in a chair drinking orange juice with ip adapter callback" />
<figcaption class="mt-2 text-center text-sm text-gray-500">with IPAdapterScaleCutoffCallback</figcaption>
</div>
</div>
## Display image after each generation step
> [!TIP]

View File

@@ -461,12 +461,12 @@ Chain it to an upscaler pipeline to increase the image resolution:
from diffusers import StableDiffusionLatentUpscalePipeline
upscaler = StableDiffusionLatentUpscalePipeline.from_pretrained(
"stabilityai/sd-x2-latent-upscaler", torch_dtype=torch.float16, use_safetensors=True
"stabilityai/sd-x2-latent-upscaler", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
)
upscaler.enable_model_cpu_offload()
upscaler.enable_xformers_memory_efficient_attention()
image_2 = upscaler(prompt, image=image_1).images[0]
image_2 = upscaler(prompt, image=image_1, output_type="latent").images[0]
```
Finally, chain it to a super-resolution pipeline to further enhance the resolution:

View File

@@ -1,6 +1,4 @@
<!--
Copyright 2023-2025 Marigold Team, ETH Zürich. All rights reserved.
Copyright 2024-2025 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 Marigold authors and The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,38 +10,31 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Marigold Computer Vision
# Marigold Pipelines for Computer Vision Tasks
**Marigold** is a diffusion-based [method](https://huggingface.co/papers/2312.02145) and a collection of [pipelines](../api/pipelines/marigold) designed for
dense computer vision tasks, including **monocular depth prediction**, **surface normals estimation**, and **intrinsic
image decomposition**.
[Marigold](../api/pipelines/marigold) is a novel diffusion-based dense prediction approach, and a set of pipelines for various computer vision tasks, such as monocular depth estimation.
This guide will walk you through using Marigold to generate fast and high-quality predictions for images and videos.
This guide will show you how to use Marigold to obtain fast and high-quality predictions for images and videos.
Each pipeline is tailored for a specific computer vision task, processing an input RGB image and generating a
corresponding prediction.
Currently, the following computer vision tasks are implemented:
Each pipeline supports one Computer Vision task, which takes an input RGB image as input and produces a *prediction* of the modality of interest, such as a depth map of the input image.
Currently, the following tasks are implemented:
| Pipeline | Recommended Model Checkpoints | Spaces (Interactive Apps) | Predicted Modalities |
|---------------------------------------------------------------------------------------------------------------------------------------------------|---------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|:------------------------------------------------------------------------------------:|------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [prs-eth/marigold-depth-v1-1](https://huggingface.co/prs-eth/marigold-depth-v1-1) | [Depth Estimation](https://huggingface.co/spaces/prs-eth/marigold) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) |
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [prs-eth/marigold-normals-v1-1](https://huggingface.co/prs-eth/marigold-normals-v1-1) | [Surface Normals Estimation](https://huggingface.co/spaces/prs-eth/marigold-normals) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) |
| [MarigoldIntrinsicsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_intrinsics.py) | [prs-eth/marigold-iid-appearance-v1-1](https://huggingface.co/prs-eth/marigold-iid-appearance-v1-1),<br>[prs-eth/marigold-iid-lighting-v1-1](https://huggingface.co/prs-eth/marigold-iid-lighting-v1-1) | [Intrinsic Image Decomposition](https://huggingface.co/spaces/prs-eth/marigold-iid) | [Albedo](https://en.wikipedia.org/wiki/Albedo), [Materials](https://www.n.aiq3d.com/wiki/roughnessmetalnessao-map), [Lighting](https://en.wikipedia.org/wiki/Diffuse_reflection) |
| Pipeline | Predicted Modalities | Demos |
|---------------------------------------------------------------------------------------------------------------------------------------------|------------------------------------------------------------------------------------------------------------------|:--------------------------------------------------------------------------------------------------------------------------------------------------:|
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-lcm), [Slow Original Demo (DDIM)](https://huggingface.co/spaces/prs-eth/marigold) |
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-normals-lcm) |
All original checkpoints are available under the [PRS-ETH](https://huggingface.co/prs-eth/) organization on Hugging Face.
They are designed for use with diffusers pipelines and the [original codebase](https://github.com/prs-eth/marigold), which can also be used to train
new model checkpoints.
The following is a summary of the recommended checkpoints, all of which produce reliable results with 1 to 4 steps.
The original checkpoints can be found under the [PRS-ETH](https://huggingface.co/prs-eth/) Hugging Face organization.
These checkpoints are meant to work with diffusers pipelines and the [original codebase](https://github.com/prs-eth/marigold).
The original code can also be used to train new checkpoints.
| Checkpoint | Modality | Comment |
|-----------------------------------------------------------------------------------------------------|--------------|-------------------------------------------------------------------------------------------------------------------------------------------------------------------|
| [prs-eth/marigold-depth-v1-1](https://huggingface.co/prs-eth/marigold-depth-v1-1) | Depth | Affine-invariant depth prediction assigns each pixel a value between 0 (near plane) and 1 (far plane), with both planes determined by the model during inference. |
| [prs-eth/marigold-normals-v0-1](https://huggingface.co/prs-eth/marigold-normals-v0-1) | Normals | The surface normals predictions are unit-length 3D vectors in the screen space camera, with values in the range from -1 to 1. |
| [prs-eth/marigold-iid-appearance-v1-1](https://huggingface.co/prs-eth/marigold-iid-appearance-v1-1) | Intrinsics | InteriorVerse decomposition is comprised of Albedo and two BRDF material properties: Roughness and Metallicity. |
| [prs-eth/marigold-iid-lighting-v1-1](https://huggingface.co/prs-eth/marigold-iid-lighting-v1-1) | Intrinsics | HyperSim decomposition of an image \\(I\\) is comprised of Albedo \\(A\\), Diffuse shading \\(S\\), and Non-diffuse residual \\(R\\): \\(I = A*S+R\\). |
The examples below are mostly given for depth prediction, but they can be universally applied to other supported
modalities.
| Checkpoint | Modality | Comment |
|-----------------------------------------------------------------------------------------------|----------|--------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
| [prs-eth/marigold-v1-0](https://huggingface.co/prs-eth/marigold-v1-0) | Depth | The first Marigold Depth checkpoint, which predicts *affine-invariant depth* maps. The performance of this checkpoint in benchmarks was studied in the original [paper](https://huggingface.co/papers/2312.02145). Designed to be used with the `DDIMScheduler` at inference, it requires at least 10 steps to get reliable predictions. Affine-invariant depth prediction has a range of values in each pixel between 0 (near plane) and 1 (far plane); both planes are chosen by the model as part of the inference process. See the `MarigoldImageProcessor` reference for visualization utilities. |
| [prs-eth/marigold-depth-lcm-v1-0](https://huggingface.co/prs-eth/marigold-depth-lcm-v1-0) | Depth | The fast Marigold Depth checkpoint, fine-tuned from `prs-eth/marigold-v1-0`. Designed to be used with the `LCMScheduler` at inference, it requires as little as 1 step to get reliable predictions. The prediction reliability saturates at 4 steps and declines after that. |
| [prs-eth/marigold-normals-v0-1](https://huggingface.co/prs-eth/marigold-normals-v0-1) | Normals | A preview checkpoint for the Marigold Normals pipeline. Designed to be used with the `DDIMScheduler` at inference, it requires at least 10 steps to get reliable predictions. The surface normals predictions are unit-length 3D vectors with values in the range from -1 to 1. *This checkpoint will be phased out after the release of `v1-0` version.* |
| [prs-eth/marigold-normals-lcm-v0-1](https://huggingface.co/prs-eth/marigold-normals-lcm-v0-1) | Normals | The fast Marigold Normals checkpoint, fine-tuned from `prs-eth/marigold-normals-v0-1`. Designed to be used with the `LCMScheduler` at inference, it requires as little as 1 step to get reliable predictions. The prediction reliability saturates at 4 steps and declines after that. *This checkpoint will be phased out after the release of `v1-0` version.* |
The examples below are mostly given for depth prediction, but they can be universally applied with other supported modalities.
We showcase the predictions using the same input image of Albert Einstein generated by Midjourney.
This makes it easier to compare visualizations of the predictions across various modalities and checkpoints.
@@ -56,21 +47,19 @@ This makes it easier to compare visualizations of the predictions across various
</div>
</div>
## Depth Prediction
### Depth Prediction Quick Start
To get a depth prediction, load the `prs-eth/marigold-depth-v1-1` checkpoint into [`MarigoldDepthPipeline`],
put the image through the pipeline, and save the predictions:
To get the first depth prediction, load `prs-eth/marigold-depth-lcm-v1-0` checkpoint into `MarigoldDepthPipeline` pipeline, put the image through the pipeline, and save the predictions:
```python
import diffusers
import torch
pipe = diffusers.MarigoldDepthPipeline.from_pretrained(
"prs-eth/marigold-depth-v1-1", variant="fp16", torch_dtype=torch.float16
"prs-eth/marigold-depth-lcm-v1-0", variant="fp16", torch_dtype=torch.float16
).to("cuda")
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
depth = pipe(image)
vis = pipe.image_processor.visualize_depth(depth.prediction)
@@ -80,13 +69,10 @@ depth_16bit = pipe.image_processor.export_depth_to_16bit_png(depth.prediction)
depth_16bit[0].save("einstein_depth_16bit.png")
```
The [`~pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_depth`] function applies one of
[matplotlib's colormaps](https://matplotlib.org/stable/users/explain/colors/colormaps.html) (`Spectral` by default) to map the predicted pixel values from a single-channel `[0, 1]`
depth range into an RGB image.
With the `Spectral` colormap, pixels with near depth are painted red, and far pixels are blue.
The visualization function for depth [`~pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_depth`] applies one of [matplotlib's colormaps](https://matplotlib.org/stable/users/explain/colors/colormaps.html) (`Spectral` by default) to map the predicted pixel values from a single-channel `[0, 1]` depth range into an RGB image.
With the `Spectral` colormap, pixels with near depth are painted red, and far pixels are assigned blue color.
The 16-bit PNG file stores the single channel values mapped linearly from the `[0, 1]` range into `[0, 65535]`.
Below are the raw and the visualized predictions. The darker and closer areas (mustache) are easier to distinguish in
the visualization.
Below are the raw and the visualized predictions; as can be seen, dark areas (mustache) are easier to distinguish in the visualization:
<div class="flex gap-4">
<div style="flex: 1 1 50%; max-width: 50%;">
@@ -103,33 +89,28 @@ the visualization.
</div>
</div>
## Surface Normals Estimation
### Surface Normals Prediction Quick Start
Load the `prs-eth/marigold-normals-v1-1` checkpoint into [`MarigoldNormalsPipeline`], put the image through the
pipeline, and save the predictions:
Load `prs-eth/marigold-normals-lcm-v0-1` checkpoint into `MarigoldNormalsPipeline` pipeline, put the image through the pipeline, and save the predictions:
```python
import diffusers
import torch
pipe = diffusers.MarigoldNormalsPipeline.from_pretrained(
"prs-eth/marigold-normals-v1-1", variant="fp16", torch_dtype=torch.float16
"prs-eth/marigold-normals-lcm-v0-1", variant="fp16", torch_dtype=torch.float16
).to("cuda")
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
normals = pipe(image)
vis = pipe.image_processor.visualize_normals(normals.prediction)
vis[0].save("einstein_normals.png")
```
The [`~pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_normals`] maps the three-dimensional
prediction with pixel values in the range `[-1, 1]` into an RGB image.
The visualization function supports flipping surface normals axes to make the visualization compatible with other
choices of the frame of reference.
Conceptually, each pixel is painted according to the surface normal vector in the frame of reference, where `X` axis
points right, `Y` axis points up, and `Z` axis points at the viewer.
The visualization function for normals [`~pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_normals`] maps the three-dimensional prediction with pixel values in the range `[-1, 1]` into an RGB image.
The visualization function supports flipping surface normals axes to make the visualization compatible with other choices of the frame of reference.
Conceptually, each pixel is painted according to the surface normal vector in the frame of reference, where `X` axis points right, `Y` axis points up, and `Z` axis points at the viewer.
Below is the visualized prediction:
<div class="flex gap-4" style="justify-content: center; width: 100%;">
@@ -141,121 +122,25 @@ Below is the visualized prediction:
</div>
</div>
In this example, the nose tip almost certainly has a point on the surface, in which the surface normal vector points
straight at the viewer, meaning that its coordinates are `[0, 0, 1]`.
In this example, the nose tip almost certainly has a point on the surface, in which the surface normal vector points straight at the viewer, meaning that its coordinates are `[0, 0, 1]`.
This vector maps to the RGB `[128, 128, 255]`, which corresponds to the violet-blue color.
Similarly, a surface normal on the cheek in the right part of the image has a large `X` component, which increases the
red hue.
Similarly, a surface normal on the cheek in the right part of the image has a large `X` component, which increases the red hue.
Points on the shoulders pointing up with a large `Y` promote green color.
## Intrinsic Image Decomposition
### Speeding up inference
Marigold provides two models for Intrinsic Image Decomposition (IID): "Appearance" and "Lighting".
Each model produces Albedo maps, derived from InteriorVerse and Hypersim annotations, respectively.
- The "Appearance" model also estimates Material properties: Roughness and Metallicity.
- The "Lighting" model generates Diffuse Shading and Non-diffuse Residual.
Here is the sample code saving predictions made by the "Appearance" model:
```python
import diffusers
import torch
pipe = diffusers.MarigoldIntrinsicsPipeline.from_pretrained(
"prs-eth/marigold-iid-appearance-v1-1", variant="fp16", torch_dtype=torch.float16
).to("cuda")
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
intrinsics = pipe(image)
vis = pipe.image_processor.visualize_intrinsics(intrinsics.prediction, pipe.target_properties)
vis[0]["albedo"].save("einstein_albedo.png")
vis[0]["roughness"].save("einstein_roughness.png")
vis[0]["metallicity"].save("einstein_metallicity.png")
```
Another example demonstrating the predictions made by the "Lighting" model:
```python
import diffusers
import torch
pipe = diffusers.MarigoldIntrinsicsPipeline.from_pretrained(
"prs-eth/marigold-iid-lighting-v1-1", variant="fp16", torch_dtype=torch.float16
).to("cuda")
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
intrinsics = pipe(image)
vis = pipe.image_processor.visualize_intrinsics(intrinsics.prediction, pipe.target_properties)
vis[0]["albedo"].save("einstein_albedo.png")
vis[0]["shading"].save("einstein_shading.png")
vis[0]["residual"].save("einstein_residual.png")
```
Both models share the same pipeline while supporting different decomposition types.
The exact decomposition parameterization (e.g., sRGB vs. linear space) is stored in the
`pipe.target_properties` dictionary, which is passed into the
[`~pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_intrinsics`] function.
Below are some examples showcasing the predicted decomposition outputs.
All modalities can be inspected in the
[Intrinsic Image Decomposition](https://huggingface.co/spaces/prs-eth/marigold-iid) Space.
<div class="flex gap-4">
<div style="flex: 1 1 50%; max-width: 50%;">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/8c7986eaaab5eb9604eb88336311f46a7b0ff5ab/marigold/marigold_einstein_albedo.png"/>
<figcaption class="mt-1 text-center text-sm text-gray-500">
Predicted albedo ("Appearance" model)
</figcaption>
</div>
<div style="flex: 1 1 50%; max-width: 50%;">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/8c7986eaaab5eb9604eb88336311f46a7b0ff5ab/marigold/marigold_einstein_diffuse.png"/>
<figcaption class="mt-1 text-center text-sm text-gray-500">
Predicted diffuse shading ("Lighting" model)
</figcaption>
</div>
</div>
## Speeding up inference
The above quick start snippets are already optimized for quality and speed, loading the checkpoint, utilizing the
`fp16` variant of weights and computation, and performing the default number (4) of denoising diffusion steps.
The first step to accelerate inference, at the expense of prediction quality, is to reduce the denoising diffusion
steps to the minimum:
The above quick start snippets are already optimized for speed: they load the LCM checkpoint, use the `fp16` variant of weights and computation, and perform just one denoising diffusion step.
The `pipe(image)` call completes in 280ms on RTX 3090 GPU.
Internally, the input image is encoded with the Stable Diffusion VAE encoder, then the U-Net performs one denoising step, and finally, the prediction latent is decoded with the VAE decoder into pixel space.
In this case, two out of three module calls are dedicated to converting between pixel and latent space of LDM.
Because Marigold's latent space is compatible with the base Stable Diffusion, it is possible to speed up the pipeline call by more than 3x (85ms on RTX 3090) by using a [lightweight replacement of the SD VAE](../api/models/autoencoder_tiny):
```diff
import diffusers
import torch
pipe = diffusers.MarigoldDepthPipeline.from_pretrained(
"prs-eth/marigold-depth-v1-1", variant="fp16", torch_dtype=torch.float16
).to("cuda")
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
- depth = pipe(image)
+ depth = pipe(image, num_inference_steps=1)
```
With this change, the `pipe` call completes in 280ms on RTX 3090 GPU.
Internally, the input image is first encoded using the Stable Diffusion VAE encoder, followed by a single denoising
step performed by the U-Net.
Finally, the prediction latent is decoded with the VAE decoder into pixel space.
In this setup, two out of three module calls are dedicated to converting between the pixel and latent spaces of the LDM.
Since Marigold's latent space is compatible with Stable Diffusion 2.0, inference can be accelerated by more than 3x,
reducing the call time to 85ms on an RTX 3090, by using a [lightweight replacement of the SD VAE](../api/models/autoencoder_tiny).
Note that using a lightweight VAE may slightly reduce the visual quality of the predictions.
```diff
import diffusers
import torch
pipe = diffusers.MarigoldDepthPipeline.from_pretrained(
"prs-eth/marigold-depth-v1-1", variant="fp16", torch_dtype=torch.float16
"prs-eth/marigold-depth-lcm-v1-0", variant="fp16", torch_dtype=torch.float16
).to("cuda")
+ pipe.vae = diffusers.AutoencoderTiny.from_pretrained(
@@ -263,77 +148,78 @@ Note that using a lightweight VAE may slightly reduce the visual quality of the
+ ).cuda()
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
depth = pipe(image, num_inference_steps=1)
depth = pipe(image)
```
So far, we have optimized the number of diffusion steps and model components. Self-attention operations account for a
significant portion of computations.
Speeding them up can be achieved by using a more efficient attention processor:
As suggested in [Optimizations](../optimization/torch2.0#torch.compile), adding `torch.compile` may squeeze extra performance depending on the target hardware:
```diff
import diffusers
import torch
+ from diffusers.models.attention_processor import AttnProcessor2_0
pipe = diffusers.MarigoldDepthPipeline.from_pretrained(
"prs-eth/marigold-depth-v1-1", variant="fp16", torch_dtype=torch.float16
"prs-eth/marigold-depth-lcm-v1-0", variant="fp16", torch_dtype=torch.float16
).to("cuda")
+ pipe.vae.set_attn_processor(AttnProcessor2_0())
+ pipe.unet.set_attn_processor(AttnProcessor2_0())
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
depth = pipe(image, num_inference_steps=1)
```
Finally, as suggested in [Optimizations](../optimization/torch2.0#torch.compile), enabling `torch.compile` can further enhance performance depending on
the target hardware.
However, compilation incurs a significant overhead during the first pipeline invocation, making it beneficial only when
the same pipeline instance is called repeatedly, such as within a loop.
```diff
import diffusers
import torch
from diffusers.models.attention_processor import AttnProcessor2_0
pipe = diffusers.MarigoldDepthPipeline.from_pretrained(
"prs-eth/marigold-depth-v1-1", variant="fp16", torch_dtype=torch.float16
).to("cuda")
pipe.vae.set_attn_processor(AttnProcessor2_0())
pipe.unet.set_attn_processor(AttnProcessor2_0())
+ pipe.vae = torch.compile(pipe.vae, mode="reduce-overhead", fullgraph=True)
+ pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
depth = pipe(image, num_inference_steps=1)
depth = pipe(image)
```
## Qualitative Comparison with Depth Anything
With the above speed optimizations, Marigold delivers predictions with more details and faster than [Depth Anything](https://huggingface.co/docs/transformers/main/en/model_doc/depth_anything) with the largest checkpoint [LiheYoung/depth-anything-large-hf](https://huggingface.co/LiheYoung/depth-anything-large-hf):
<div class="flex gap-4">
<div style="flex: 1 1 50%; max-width: 50%;">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/marigold/marigold_einstein_lcm_depth.png"/>
<figcaption class="mt-1 text-center text-sm text-gray-500">
Marigold LCM fp16 with Tiny AutoEncoder
</figcaption>
</div>
<div style="flex: 1 1 50%; max-width: 50%;">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/marigold/einstein_depthanything_large.png"/>
<figcaption class="mt-1 text-center text-sm text-gray-500">
Depth Anything Large
</figcaption>
</div>
</div>
## Maximizing Precision and Ensembling
Marigold pipelines have a built-in ensembling mechanism combining multiple predictions from different random latents.
This is a brute-force way of improving the precision of predictions, capitalizing on the generative nature of diffusion.
The ensembling path is activated automatically when the `ensemble_size` argument is set greater or equal than `3`.
The ensembling path is activated automatically when the `ensemble_size` argument is set greater than `1`.
When aiming for maximum precision, it makes sense to adjust `num_inference_steps` simultaneously with `ensemble_size`.
The recommended values vary across checkpoints but primarily depend on the scheduler type.
The effect of ensembling is particularly well-seen with surface normals:
```diff
import diffusers
```python
import diffusers
pipe = diffusers.MarigoldNormalsPipeline.from_pretrained("prs-eth/marigold-normals-v1-1").to("cuda")
model_path = "prs-eth/marigold-normals-v1-0"
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
model_paper_kwargs = {
diffusers.schedulers.DDIMScheduler: {
"num_inference_steps": 10,
"ensemble_size": 10,
},
diffusers.schedulers.LCMScheduler: {
"num_inference_steps": 4,
"ensemble_size": 5,
},
}
- depth = pipe(image)
+ depth = pipe(image, num_inference_steps=10, ensemble_size=5)
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
vis = pipe.image_processor.visualize_normals(depth.prediction)
vis[0].save("einstein_normals.png")
pipe = diffusers.MarigoldNormalsPipeline.from_pretrained(model_path).to("cuda")
pipe_kwargs = model_paper_kwargs[type(pipe.scheduler)]
depth = pipe(image, **pipe_kwargs)
vis = pipe.image_processor.visualize_normals(depth.prediction)
vis[0].save("einstein_normals.png")
```
<div class="flex gap-4">
@@ -351,16 +237,93 @@ The effect of ensembling is particularly well-seen with surface normals:
</div>
</div>
As can be seen, all areas with fine-grained structurers, such as hair, got more conservative and on average more
correct predictions.
As can be seen, all areas with fine-grained structurers, such as hair, got more conservative and on average more correct predictions.
Such a result is more suitable for precision-sensitive downstream tasks, such as 3D reconstruction.
## Quantitative Evaluation
To evaluate Marigold quantitatively in standard leaderboards and benchmarks (such as NYU, KITTI, and other datasets), follow the evaluation protocol outlined in the paper: load the full precision fp32 model and use appropriate values for `num_inference_steps` and `ensemble_size`.
Optionally seed randomness to ensure reproducibility. Maximizing `batch_size` will deliver maximum device utilization.
```python
import diffusers
import torch
device = "cuda"
seed = 2024
model_path = "prs-eth/marigold-v1-0"
model_paper_kwargs = {
diffusers.schedulers.DDIMScheduler: {
"num_inference_steps": 50,
"ensemble_size": 10,
},
diffusers.schedulers.LCMScheduler: {
"num_inference_steps": 4,
"ensemble_size": 10,
},
}
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
generator = torch.Generator(device=device).manual_seed(seed)
pipe = diffusers.MarigoldDepthPipeline.from_pretrained(model_path).to(device)
pipe_kwargs = model_paper_kwargs[type(pipe.scheduler)]
depth = pipe(image, generator=generator, **pipe_kwargs)
# evaluate metrics
```
## Using Predictive Uncertainty
The ensembling mechanism built into Marigold pipelines combines multiple predictions obtained from different random latents.
As a side effect, it can be used to quantify epistemic (model) uncertainty; simply specify `ensemble_size` greater than 1 and set `output_uncertainty=True`.
The resulting uncertainty will be available in the `uncertainty` field of the output.
It can be visualized as follows:
```python
import diffusers
import torch
pipe = diffusers.MarigoldDepthPipeline.from_pretrained(
"prs-eth/marigold-depth-lcm-v1-0", variant="fp16", torch_dtype=torch.float16
).to("cuda")
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
depth = pipe(
image,
ensemble_size=10, # any number greater than 1; higher values yield higher precision
output_uncertainty=True,
)
uncertainty = pipe.image_processor.visualize_uncertainty(depth.uncertainty)
uncertainty[0].save("einstein_depth_uncertainty.png")
```
<div class="flex gap-4">
<div style="flex: 1 1 50%; max-width: 50%;">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/marigold/marigold_einstein_depth_uncertainty.png"/>
<figcaption class="mt-1 text-center text-sm text-gray-500">
Depth uncertainty
</figcaption>
</div>
<div style="flex: 1 1 50%; max-width: 50%;">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/marigold/marigold_einstein_normals_uncertainty.png"/>
<figcaption class="mt-1 text-center text-sm text-gray-500">
Surface normals uncertainty
</figcaption>
</div>
</div>
The interpretation of uncertainty is easy: higher values (white) correspond to pixels, where the model struggles to make consistent predictions.
Evidently, the depth model is the least confident around edges with discontinuity, where the object depth changes drastically.
The surface normals model is the least confident in fine-grained structures, such as hair, and dark areas, such as the collar.
## Frame-by-frame Video Processing with Temporal Consistency
Due to Marigold's generative nature, each prediction is unique and defined by the random noise sampled for the latent
initialization.
This becomes an obvious drawback compared to traditional end-to-end dense regression networks, as exemplified in the
following videos:
Due to Marigold's generative nature, each prediction is unique and defined by the random noise sampled for the latent initialization.
This becomes an obvious drawback compared to traditional end-to-end dense regression networks, as exemplified in the following videos:
<div class="flex gap-4">
<div style="flex: 1 1 50%; max-width: 50%;">
@@ -373,32 +336,26 @@ following videos:
</div>
</div>
To address this issue, it is possible to pass `latents` argument to the pipelines, which defines the starting point of
diffusion.
Empirically, we found that a convex combination of the very same starting point noise latent and the latent
corresponding to the previous frame prediction give sufficiently smooth results, as implemented in the snippet below:
To address this issue, it is possible to pass `latents` argument to the pipelines, which defines the starting point of diffusion.
Empirically, we found that a convex combination of the very same starting point noise latent and the latent corresponding to the previous frame prediction give sufficiently smooth results, as implemented in the snippet below:
```python
import imageio
import diffusers
import torch
from diffusers.models.attention_processor import AttnProcessor2_0
from PIL import Image
from tqdm import tqdm
import diffusers
import torch
device = "cuda"
path_in = "https://huggingface.co/spaces/prs-eth/marigold-lcm/resolve/c7adb5427947d2680944f898cd91d386bf0d4924/files/video/obama.mp4"
path_in = "obama.mp4"
path_out = "obama_depth.gif"
pipe = diffusers.MarigoldDepthPipeline.from_pretrained(
"prs-eth/marigold-depth-v1-1", variant="fp16", torch_dtype=torch.float16
"prs-eth/marigold-depth-lcm-v1-0", variant="fp16", torch_dtype=torch.float16
).to(device)
pipe.vae = diffusers.AutoencoderTiny.from_pretrained(
"madebyollin/taesd", torch_dtype=torch.float16
).to(device)
pipe.unet.set_attn_processor(AttnProcessor2_0())
pipe.vae = torch.compile(pipe.vae, mode="reduce-overhead", fullgraph=True)
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
pipe.set_progress_bar_config(disable=True)
with imageio.get_reader(path_in) as reader:
@@ -416,11 +373,7 @@ with imageio.get_reader(path_in) as reader:
latents = 0.9 * latents + 0.1 * last_frame_latent
depth = pipe(
frame,
num_inference_steps=1,
match_input_resolution=False,
latents=latents,
output_latent=True,
frame, match_input_resolution=False, latents=latents, output_latent=True
)
last_frame_latent = depth.latent
out.append(pipe.image_processor.visualize_depth(depth.prediction)[0])
@@ -429,8 +382,7 @@ with imageio.get_reader(path_in) as reader:
```
Here, the diffusion process starts from the given computed latent.
The pipeline sets `output_latent=True` to access `out.latent` and computes its contribution to the next frame's latent
initialization.
The pipeline sets `output_latent=True` to access `out.latent` and computes its contribution to the next frame's latent initialization.
The result is much more stable now:
<div class="flex gap-4">
@@ -462,7 +414,7 @@ image = diffusers.utils.load_image(
)
pipe = diffusers.MarigoldDepthPipeline.from_pretrained(
"prs-eth/marigold-depth-v1-1", torch_dtype=torch.float16, variant="fp16"
"prs-eth/marigold-depth-lcm-v1-0", torch_dtype=torch.float16, variant="fp16"
).to(device)
depth_image = pipe(image, generator=generator).prediction
@@ -511,95 +463,4 @@ controlnet_out[0].save("motorcycle_controlnet_out.png")
</div>
</div>
## Quantitative Evaluation
To evaluate Marigold quantitatively in standard leaderboards and benchmarks (such as NYU, KITTI, and other datasets),
follow the evaluation protocol outlined in the paper: load the full precision fp32 model and use appropriate values
for `num_inference_steps` and `ensemble_size`.
Optionally seed randomness to ensure reproducibility.
Maximizing `batch_size` will deliver maximum device utilization.
```python
import diffusers
import torch
device = "cuda"
seed = 2024
generator = torch.Generator(device=device).manual_seed(seed)
pipe = diffusers.MarigoldDepthPipeline.from_pretrained("prs-eth/marigold-depth-v1-1").to(device)
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
depth = pipe(
image,
num_inference_steps=4, # set according to the evaluation protocol from the paper
ensemble_size=10, # set according to the evaluation protocol from the paper
generator=generator,
)
# evaluate metrics
```
## Using Predictive Uncertainty
The ensembling mechanism built into Marigold pipelines combines multiple predictions obtained from different random
latents.
As a side effect, it can be used to quantify epistemic (model) uncertainty; simply specify `ensemble_size` greater
or equal than 3 and set `output_uncertainty=True`.
The resulting uncertainty will be available in the `uncertainty` field of the output.
It can be visualized as follows:
```python
import diffusers
import torch
pipe = diffusers.MarigoldDepthPipeline.from_pretrained(
"prs-eth/marigold-depth-v1-1", variant="fp16", torch_dtype=torch.float16
).to("cuda")
image = diffusers.utils.load_image("https://marigoldmonodepth.github.io/images/einstein.jpg")
depth = pipe(
image,
ensemble_size=10, # any number >= 3
output_uncertainty=True,
)
uncertainty = pipe.image_processor.visualize_uncertainty(depth.uncertainty)
uncertainty[0].save("einstein_depth_uncertainty.png")
```
<div class="flex gap-4">
<div style="flex: 1 1 33%; max-width: 33%;">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/marigold/marigold_einstein_depth_uncertainty.png"/>
<figcaption class="mt-1 text-center text-sm text-gray-500">
Depth uncertainty
</figcaption>
</div>
<div style="flex: 1 1 33%; max-width: 33%;">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/marigold/marigold_einstein_normals_uncertainty.png"/>
<figcaption class="mt-1 text-center text-sm text-gray-500">
Surface normals uncertainty
</figcaption>
</div>
<div style="flex: 1 1 33%; max-width: 33%;">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/4f83035d84a24e5ec44fdda129b1d51eba12ce04/marigold/marigold_einstein_albedo_uncertainty.png"/>
<figcaption class="mt-1 text-center text-sm text-gray-500">
Albedo uncertainty
</figcaption>
</div>
</div>
The interpretation of uncertainty is easy: higher values (white) correspond to pixels, where the model struggles to
make consistent predictions.
- The depth model exhibits the most uncertainty around discontinuities, where object depth changes abruptly.
- The surface normals model is least confident in fine-grained structures like hair and in dark regions such as the
collar area.
- Albedo uncertainty is represented as an RGB image, as it captures uncertainty independently for each color channel,
unlike depth and surface normals. It is also higher in shaded regions and at discontinuities.
## Conclusion
We hope Marigold proves valuable for your downstream tasks, whether as part of a broader generative workflow or for
perception-based applications like 3D reconstruction.
Hopefully, you will find Marigold useful for solving your downstream tasks, be it a part of a more broad generative workflow, or a perception task, such as 3D reconstruction.

View File

@@ -1,317 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# OmniGen
OmniGen is an image generation model. Unlike existing text-to-image models, OmniGen is a single model designed to handle a variety of tasks (e.g., text-to-image, image editing, controllable generation). It has the following features:
- Minimalist model architecture, consisting of only a VAE and a transformer module, for joint modeling of text and images.
- Support for multimodal inputs. It can process any text-image mixed data as instructions for image generation, rather than relying solely on text.
For more information, please refer to the [paper](https://arxiv.org/pdf/2409.11340).
This guide will walk you through using OmniGen for various tasks and use cases.
## Load model checkpoints
Model weights may be stored in separate subfolders on the Hub or locally, in which case, you should use the [`~DiffusionPipeline.from_pretrained`] method.
```python
import torch
from diffusers import OmniGenPipeline
pipe = OmniGenPipeline.from_pretrained("Shitao/OmniGen-v1-diffusers", torch_dtype=torch.bfloat16)
```
## Text-to-image
For text-to-image, pass a text prompt. By default, OmniGen generates a 1024x1024 image.
You can try setting the `height` and `width` parameters to generate images with different size.
```python
import torch
from diffusers import OmniGenPipeline
pipe = OmniGenPipeline.from_pretrained(
"Shitao/OmniGen-v1-diffusers",
torch_dtype=torch.bfloat16
)
pipe.to("cuda")
prompt = "Realistic photo. A young woman sits on a sofa, holding a book and facing the camera. She wears delicate silver hoop earrings adorned with tiny, sparkling diamonds that catch the light, with her long chestnut hair cascading over her shoulders. Her eyes are focused and gentle, framed by long, dark lashes. She is dressed in a cozy cream sweater, which complements her warm, inviting smile. Behind her, there is a table with a cup of water in a sleek, minimalist blue mug. The background is a serene indoor setting with soft natural light filtering through a window, adorned with tasteful art and flowers, creating a cozy and peaceful ambiance. 4K, HD."
image = pipe(
prompt=prompt,
height=1024,
width=1024,
guidance_scale=3,
generator=torch.Generator(device="cpu").manual_seed(111),
).images[0]
image.save("output.png")
```
<div class="flex justify-center">
<img src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png" alt="generated image"/>
</div>
## Image edit
OmniGen supports multimodal inputs.
When the input includes an image, you need to add a placeholder `<img><|image_1|></img>` in the text prompt to represent the image.
It is recommended to enable `use_input_image_size_as_output` to keep the edited image the same size as the original image.
```python
import torch
from diffusers import OmniGenPipeline
from diffusers.utils import load_image
pipe = OmniGenPipeline.from_pretrained(
"Shitao/OmniGen-v1-diffusers",
torch_dtype=torch.bfloat16
)
pipe.to("cuda")
prompt="<img><|image_1|></img> Remove the woman's earrings. Replace the mug with a clear glass filled with sparkling iced cola."
input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png")]
image = pipe(
prompt=prompt,
input_images=input_images,
guidance_scale=2,
img_guidance_scale=1.6,
use_input_image_size_as_output=True,
generator=torch.Generator(device="cpu").manual_seed(222)
).images[0]
image.save("output.png")
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">original image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/edit.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">edited image</figcaption>
</div>
</div>
OmniGen has some interesting features, such as visual reasoning, as shown in the example below.
```python
prompt="If the woman is thirsty, what should she take? Find it in the image and highlight it in blue. <img><|image_1|></img>"
input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/edit.png")]
image = pipe(
prompt=prompt,
input_images=input_images,
guidance_scale=2,
img_guidance_scale=1.6,
use_input_image_size_as_output=True,
generator=torch.Generator(device="cpu").manual_seed(0)
).images[0]
image.save("output.png")
```
<div class="flex justify-center">
<img src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/reasoning.png" alt="generated image"/>
</div>
## Controllable generation
OmniGen can handle several classic computer vision tasks. As shown below, OmniGen can detect human skeletons in input images, which can be used as control conditions to generate new images.
```python
import torch
from diffusers import OmniGenPipeline
from diffusers.utils import load_image
pipe = OmniGenPipeline.from_pretrained(
"Shitao/OmniGen-v1-diffusers",
torch_dtype=torch.bfloat16
)
pipe.to("cuda")
prompt="Detect the skeleton of human in this image: <img><|image_1|></img>"
input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/edit.png")]
image1 = pipe(
prompt=prompt,
input_images=input_images,
guidance_scale=2,
img_guidance_scale=1.6,
use_input_image_size_as_output=True,
generator=torch.Generator(device="cpu").manual_seed(333)
).images[0]
image1.save("image1.png")
prompt="Generate a new photo using the following picture and text as conditions: <img><|image_1|></img>\n A young boy is sitting on a sofa in the library, holding a book. His hair is neatly combed, and a faint smile plays on his lips, with a few freckles scattered across his cheeks. The library is quiet, with rows of shelves filled with books stretching out behind him."
input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/skeletal.png")]
image2 = pipe(
prompt=prompt,
input_images=input_images,
guidance_scale=2,
img_guidance_scale=1.6,
use_input_image_size_as_output=True,
generator=torch.Generator(device="cpu").manual_seed(333)
).images[0]
image2.save("image2.png")
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/edit.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">original image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/skeletal.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">detected skeleton</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/skeletal2img.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">skeleton to image</figcaption>
</div>
</div>
OmniGen can also directly use relevant information from input images to generate new images.
```python
import torch
from diffusers import OmniGenPipeline
from diffusers.utils import load_image
pipe = OmniGenPipeline.from_pretrained(
"Shitao/OmniGen-v1-diffusers",
torch_dtype=torch.bfloat16
)
pipe.to("cuda")
prompt="Following the pose of this image <img><|image_1|></img>, generate a new photo: A young boy is sitting on a sofa in the library, holding a book. His hair is neatly combed, and a faint smile plays on his lips, with a few freckles scattered across his cheeks. The library is quiet, with rows of shelves filled with books stretching out behind him."
input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/edit.png")]
image = pipe(
prompt=prompt,
input_images=input_images,
guidance_scale=2,
img_guidance_scale=1.6,
use_input_image_size_as_output=True,
generator=torch.Generator(device="cpu").manual_seed(0)
).images[0]
image.save("output.png")
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/same_pose.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">generated image</figcaption>
</div>
</div>
## ID and object preserving
OmniGen can generate multiple images based on the people and objects in the input image and supports inputting multiple images simultaneously.
Additionally, OmniGen can extract desired objects from an image containing multiple objects based on instructions.
```python
import torch
from diffusers import OmniGenPipeline
from diffusers.utils import load_image
pipe = OmniGenPipeline.from_pretrained(
"Shitao/OmniGen-v1-diffusers",
torch_dtype=torch.bfloat16
)
pipe.to("cuda")
prompt="A man and a woman are sitting at a classroom desk. The man is the man with yellow hair in <img><|image_1|></img>. The woman is the woman on the left of <img><|image_2|></img>"
input_image_1 = load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/3.png")
input_image_2 = load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/4.png")
input_images=[input_image_1, input_image_2]
image = pipe(
prompt=prompt,
input_images=input_images,
height=1024,
width=1024,
guidance_scale=2.5,
img_guidance_scale=1.6,
generator=torch.Generator(device="cpu").manual_seed(666)
).images[0]
image.save("output.png")
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/3.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">input_image_1</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/4.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">input_image_2</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/id2.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">generated image</figcaption>
</div>
</div>
```py
import torch
from diffusers import OmniGenPipeline
from diffusers.utils import load_image
pipe = OmniGenPipeline.from_pretrained(
"Shitao/OmniGen-v1-diffusers",
torch_dtype=torch.bfloat16
)
pipe.to("cuda")
prompt="A woman is walking down the street, wearing a white long-sleeve blouse with lace details on the sleeves, paired with a blue pleated skirt. The woman is <img><|image_1|></img>. The long-sleeve blouse and a pleated skirt are <img><|image_2|></img>."
input_image_1 = load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/emma.jpeg")
input_image_2 = load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/dress.jpg")
input_images=[input_image_1, input_image_2]
image = pipe(
prompt=prompt,
input_images=input_images,
height=1024,
width=1024,
guidance_scale=2.5,
img_guidance_scale=1.6,
generator=torch.Generator(device="cpu").manual_seed(666)
).images[0]
image.save("output.png")
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/emma.jpeg"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">person image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/dress.jpg"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">clothe image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/tryon.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">generated image</figcaption>
</div>
</div>
## Optimization when using multiple images
For text-to-image task, OmniGen requires minimal memory and time costs (9GB memory and 31s for a 1024x1024 image on A800 GPU).
However, when using input images, the computational cost increases.
Here are some guidelines to help you reduce computational costs when using multiple images. The experiments are conducted on an A800 GPU with two input images.
Like other pipelines, you can reduce memory usage by offloading the model: `pipe.enable_model_cpu_offload()` or `pipe.enable_sequential_cpu_offload() `.
In OmniGen, you can also decrease computational overhead by reducing the `max_input_image_size`.
The memory consumption for different image sizes is shown in the table below:
| Method | Memory Usage |
|---------------------------|--------------|
| max_input_image_size=1024 | 40GB |
| max_input_image_size=512 | 17GB |
| max_input_image_size=256 | 14GB |

View File

@@ -215,7 +215,7 @@ image
Prompt weighting provides a way to emphasize or de-emphasize certain parts of a prompt, allowing for more control over the generated image. A prompt can include several concepts, which gets turned into contextualized text embeddings. The embeddings are used by the model to condition its cross-attention layers to generate an image (read the Stable Diffusion [blog post](https://huggingface.co/blog/stable_diffusion) to learn more about how it works).
Prompt weighting works by increasing or decreasing the scale of the text embedding vector that corresponds to its concept in the prompt because you may not necessarily want the model to focus on all concepts equally. The easiest way to prepare the prompt embeddings is to use [Stable Diffusion Long Prompt Weighted Embedding](https://github.com/xhinker/sd_embed) (sd_embed). Once you have the prompt-weighted embeddings, you can pass them to any pipeline that has a [prompt_embeds](https://huggingface.co/docs/diffusers/en/api/pipelines/stable_diffusion/text2img#diffusers.StableDiffusionPipeline.__call__.prompt_embeds) (and optionally [negative_prompt_embeds](https://huggingface.co/docs/diffusers/en/api/pipelines/stable_diffusion/text2img#diffusers.StableDiffusionPipeline.__call__.negative_prompt_embeds)) parameter, such as [`StableDiffusionPipeline`], [`StableDiffusionControlNetPipeline`], and [`StableDiffusionXLPipeline`].
Prompt weighting works by increasing or decreasing the scale of the text embedding vector that corresponds to its concept in the prompt because you may not necessarily want the model to focus on all concepts equally. The easiest way to prepare the prompt-weighted embeddings is to use [Compel](https://github.com/damian0815/compel), a text prompt-weighting and blending library. Once you have the prompt-weighted embeddings, you can pass them to any pipeline that has a [`prompt_embeds`](https://huggingface.co/docs/diffusers/en/api/pipelines/stable_diffusion/text2img#diffusers.StableDiffusionPipeline.__call__.prompt_embeds) (and optionally [`negative_prompt_embeds`](https://huggingface.co/docs/diffusers/en/api/pipelines/stable_diffusion/text2img#diffusers.StableDiffusionPipeline.__call__.negative_prompt_embeds)) parameter, such as [`StableDiffusionPipeline`], [`StableDiffusionControlNetPipeline`], and [`StableDiffusionXLPipeline`].
<Tip>
@@ -223,99 +223,136 @@ If your favorite pipeline doesn't have a `prompt_embeds` parameter, please open
</Tip>
This guide will show you how to weight your prompts with sd_embed.
This guide will show you how to weight and blend your prompts with Compel in 🤗 Diffusers.
Before you begin, make sure you have the latest version of sd_embed installed:
```bash
pip install git+https://github.com/xhinker/sd_embed.git@main
```
For this example, let's use [`StableDiffusionXLPipeline`].
Before you begin, make sure you have the latest version of Compel installed:
```py
from diffusers import StableDiffusionXLPipeline, UniPCMultistepScheduler
# uncomment to install in Colab
#!pip install compel --upgrade
```
For this guide, let's generate an image with the prompt `"a red cat playing with a ball"` using the [`StableDiffusionPipeline`]:
```py
from diffusers import StableDiffusionPipeline, UniPCMultistepScheduler
import torch
pipe = StableDiffusionXLPipeline.from_pretrained("Lykon/dreamshaper-xl-1-0", torch_dtype=torch.float16)
pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", use_safetensors=True)
pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
pipe.to("cuda")
```
To upweight or downweight a concept, surround the text with parentheses. More parentheses applies a heavier weight on the text. You can also append a numerical multiplier to the text to indicate how much you want to increase or decrease its weights by.
prompt = "a red cat playing with a ball"
| format | multiplier |
|---|---|
| `(hippo)` | increase by 1.1x |
| `((hippo))` | increase by 1.21x |
| `(hippo:1.5)` | increase by 1.5x |
| `(hippo:0.5)` | decrease by 4x |
generator = torch.Generator(device="cpu").manual_seed(33)
Create a prompt and use a combination of parentheses and numerical multipliers to upweight various text.
```py
from sd_embed.embedding_funcs import get_weighted_text_embeddings_sdxl
prompt = """A whimsical and creative image depicting a hybrid creature that is a mix of a waffle and a hippopotamus.
This imaginative creature features the distinctive, bulky body of a hippo,
but with a texture and appearance resembling a golden-brown, crispy waffle.
The creature might have elements like waffle squares across its skin and a syrup-like sheen.
It's set in a surreal environment that playfully combines a natural water habitat of a hippo with elements of a breakfast table setting,
possibly including oversized utensils or plates in the background.
The image should evoke a sense of playful absurdity and culinary fantasy.
"""
neg_prompt = """\
skin spots,acnes,skin blemishes,age spot,(ugly:1.2),(duplicate:1.2),(morbid:1.21),(mutilated:1.2),\
(tranny:1.2),mutated hands,(poorly drawn hands:1.5),blurry,(bad anatomy:1.2),(bad proportions:1.3),\
extra limbs,(disfigured:1.2),(missing arms:1.2),(extra legs:1.2),(fused fingers:1.5),\
(too many fingers:1.5),(unclear eyes:1.2),lowers,bad hands,missing fingers,extra digit,\
bad hands,missing fingers,(extra arms and legs),(worst quality:2),(low quality:2),\
(normal quality:2),lowres,((monochrome)),((grayscale))
"""
```
Use the `get_weighted_text_embeddings_sdxl` function to generate the prompt embeddings and the negative prompt embeddings. It'll also generated the pooled and negative pooled prompt embeddings since you're using the SDXL model.
> [!TIP]
> You can safely ignore the error message below about the token index length exceeding the models maximum sequence length. All your tokens will be used in the embedding process.
>
> ```
> Token indices sequence length is longer than the specified maximum sequence length for this model
> ```
```py
(
prompt_embeds,
prompt_neg_embeds,
pooled_prompt_embeds,
negative_pooled_prompt_embeds
) = get_weighted_text_embeddings_sdxl(
pipe,
prompt=prompt,
neg_prompt=neg_prompt
)
image = pipe(
prompt_embeds=prompt_embeds,
negative_prompt_embeds=prompt_neg_embeds,
pooled_prompt_embeds=pooled_prompt_embeds,
negative_pooled_prompt_embeds=negative_pooled_prompt_embeds,
num_inference_steps=30,
height=1024,
width=1024 + 512,
guidance_scale=4.0,
generator=torch.Generator("cuda").manual_seed(2)
).images[0]
image = pipe(prompt, generator=generator, num_inference_steps=20).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/sd_embed_sdxl.png"/>
<img class="rounded-xl" src="https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/compel/forest_0.png"/>
</div>
> [!TIP]
> Refer to the [sd_embed](https://github.com/xhinker/sd_embed) repository for additional details about long prompt weighting for FLUX.1, Stable Cascade, and Stable Diffusion 1.5.
### Weighting
You'll notice there is no "ball" in the image! Let's use compel to upweight the concept of "ball" in the prompt. Create a [`Compel`](https://github.com/damian0815/compel/blob/main/doc/compel.md#compel-objects) object, and pass it a tokenizer and text encoder:
```py
from compel import Compel
compel_proc = Compel(tokenizer=pipe.tokenizer, text_encoder=pipe.text_encoder)
```
compel uses `+` or `-` to increase or decrease the weight of a word in the prompt. To increase the weight of "ball":
<Tip>
`+` corresponds to the value `1.1`, `++` corresponds to `1.1^2`, and so on. Similarly, `-` corresponds to `0.9` and `--` corresponds to `0.9^2`. Feel free to experiment with adding more `+` or `-` in your prompt!
</Tip>
```py
prompt = "a red cat playing with a ball++"
```
Pass the prompt to `compel_proc` to create the new prompt embeddings which are passed to the pipeline:
```py
prompt_embeds = compel_proc(prompt)
generator = torch.manual_seed(33)
image = pipe(prompt_embeds=prompt_embeds, generator=generator, num_inference_steps=20).images[0]
image
```
<div class="flex justify-center">
<img class="rounded-xl" src="https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/compel/forest_1.png"/>
</div>
To downweight parts of the prompt, use the `-` suffix:
```py
prompt = "a red------- cat playing with a ball"
prompt_embeds = compel_proc(prompt)
generator = torch.manual_seed(33)
image = pipe(prompt_embeds=prompt_embeds, generator=generator, num_inference_steps=20).images[0]
image
```
<div class="flex justify-center">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/compel-neg.png"/>
</div>
You can even up or downweight multiple concepts in the same prompt:
```py
prompt = "a red cat++ playing with a ball----"
prompt_embeds = compel_proc(prompt)
generator = torch.manual_seed(33)
image = pipe(prompt_embeds=prompt_embeds, generator=generator, num_inference_steps=20).images[0]
image
```
<div class="flex justify-center">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/compel-pos-neg.png"/>
</div>
### Blending
You can also create a weighted *blend* of prompts by adding `.blend()` to a list of prompts and passing it some weights. Your blend may not always produce the result you expect because it breaks some assumptions about how the text encoder functions, so just have fun and experiment with it!
```py
prompt_embeds = compel_proc('("a red cat playing with a ball", "jungle").blend(0.7, 0.8)')
generator = torch.Generator(device="cuda").manual_seed(33)
image = pipe(prompt_embeds=prompt_embeds, generator=generator, num_inference_steps=20).images[0]
image
```
<div class="flex justify-center">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/compel-blend.png"/>
</div>
### Conjunction
A conjunction diffuses each prompt independently and concatenates their results by their weighted sum. Add `.and()` to the end of a list of prompts to create a conjunction:
```py
prompt_embeds = compel_proc('["a red cat", "playing with a", "ball"].and()')
generator = torch.Generator(device="cuda").manual_seed(55)
image = pipe(prompt_embeds=prompt_embeds, generator=generator, num_inference_steps=20).images[0]
image
```
<div class="flex justify-center">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/compel-conj.png"/>
</div>
### Textual inversion
@@ -326,63 +363,35 @@ Create a pipeline and use the [`~loaders.TextualInversionLoaderMixin.load_textua
```py
import torch
from diffusers import StableDiffusionPipeline
from compel import Compel, DiffusersTextualInversionManager
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
).to("cuda")
"stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16,
use_safetensors=True, variant="fp16").to("cuda")
pipe.load_textual_inversion("sd-concepts-library/midjourney-style")
```
Add the `<midjourney-style>` text to the prompt to trigger the textual inversion.
Compel provides a `DiffusersTextualInversionManager` class to simplify prompt weighting with textual inversion. Instantiate `DiffusersTextualInversionManager` and pass it to the `Compel` class:
```py
from sd_embed.embedding_funcs import get_weighted_text_embeddings_sd15
prompt = """<midjourney-style> A whimsical and creative image depicting a hybrid creature that is a mix of a waffle and a hippopotamus.
This imaginative creature features the distinctive, bulky body of a hippo,
but with a texture and appearance resembling a golden-brown, crispy waffle.
The creature might have elements like waffle squares across its skin and a syrup-like sheen.
It's set in a surreal environment that playfully combines a natural water habitat of a hippo with elements of a breakfast table setting,
possibly including oversized utensils or plates in the background.
The image should evoke a sense of playful absurdity and culinary fantasy.
"""
neg_prompt = """\
skin spots,acnes,skin blemishes,age spot,(ugly:1.2),(duplicate:1.2),(morbid:1.21),(mutilated:1.2),\
(tranny:1.2),mutated hands,(poorly drawn hands:1.5),blurry,(bad anatomy:1.2),(bad proportions:1.3),\
extra limbs,(disfigured:1.2),(missing arms:1.2),(extra legs:1.2),(fused fingers:1.5),\
(too many fingers:1.5),(unclear eyes:1.2),lowers,bad hands,missing fingers,extra digit,\
bad hands,missing fingers,(extra arms and legs),(worst quality:2),(low quality:2),\
(normal quality:2),lowres,((monochrome)),((grayscale))
"""
textual_inversion_manager = DiffusersTextualInversionManager(pipe)
compel_proc = Compel(
tokenizer=pipe.tokenizer,
text_encoder=pipe.text_encoder,
textual_inversion_manager=textual_inversion_manager)
```
Use the `get_weighted_text_embeddings_sd15` function to generate the prompt embeddings and the negative prompt embeddings.
Incorporate the concept to condition a prompt with using the `<concept>` syntax:
```py
(
prompt_embeds,
prompt_neg_embeds,
) = get_weighted_text_embeddings_sd15(
pipe,
prompt=prompt,
neg_prompt=neg_prompt
)
prompt_embeds = compel_proc('("A red cat++ playing with a ball <midjourney-style>")')
image = pipe(
prompt_embeds=prompt_embeds,
negative_prompt_embeds=prompt_neg_embeds,
height=768,
width=896,
guidance_scale=4.0,
generator=torch.Generator("cuda").manual_seed(2)
).images[0]
image = pipe(prompt_embeds=prompt_embeds).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/sd_embed_textual_inversion.png"/>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/compel-text-inversion.png"/>
</div>
### DreamBooth
@@ -392,44 +401,70 @@ image
```py
import torch
from diffusers import DiffusionPipeline, UniPCMultistepScheduler
from compel import Compel
pipe = DiffusionPipeline.from_pretrained("sd-dreambooth-library/dndcoverart-v1", torch_dtype=torch.float16).to("cuda")
pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
```
Depending on the model you use, you'll need to incorporate the model's unique identifier into your prompt. For example, the `dndcoverart-v1` model uses the identifier `dndcoverart`:
Create a `Compel` class with a tokenizer and text encoder, and pass your prompt to it. Depending on the model you use, you'll need to incorporate the model's unique identifier into your prompt. For example, the `dndcoverart-v1` model uses the identifier `dndcoverart`:
```py
from sd_embed.embedding_funcs import get_weighted_text_embeddings_sd15
prompt = """dndcoverart of A whimsical and creative image depicting a hybrid creature that is a mix of a waffle and a hippopotamus.
This imaginative creature features the distinctive, bulky body of a hippo,
but with a texture and appearance resembling a golden-brown, crispy waffle.
The creature might have elements like waffle squares across its skin and a syrup-like sheen.
It's set in a surreal environment that playfully combines a natural water habitat of a hippo with elements of a breakfast table setting,
possibly including oversized utensils or plates in the background.
The image should evoke a sense of playful absurdity and culinary fantasy.
"""
neg_prompt = """\
skin spots,acnes,skin blemishes,age spot,(ugly:1.2),(duplicate:1.2),(morbid:1.21),(mutilated:1.2),\
(tranny:1.2),mutated hands,(poorly drawn hands:1.5),blurry,(bad anatomy:1.2),(bad proportions:1.3),\
extra limbs,(disfigured:1.2),(missing arms:1.2),(extra legs:1.2),(fused fingers:1.5),\
(too many fingers:1.5),(unclear eyes:1.2),lowers,bad hands,missing fingers,extra digit,\
bad hands,missing fingers,(extra arms and legs),(worst quality:2),(low quality:2),\
(normal quality:2),lowres,((monochrome)),((grayscale))
"""
(
prompt_embeds
, prompt_neg_embeds
) = get_weighted_text_embeddings_sd15(
pipe
, prompt = prompt
, neg_prompt = neg_prompt
)
compel_proc = Compel(tokenizer=pipe.tokenizer, text_encoder=pipe.text_encoder)
prompt_embeds = compel_proc('("magazine cover of a dndcoverart dragon, high quality, intricate details, larry elmore art style").and()')
image = pipe(prompt_embeds=prompt_embeds).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/sd_embed_dreambooth.png"/>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/compel-dreambooth.png"/>
</div>
### Stable Diffusion XL
Stable Diffusion XL (SDXL) has two tokenizers and text encoders so it's usage is a bit different. To address this, you should pass both tokenizers and encoders to the `Compel` class:
```py
from compel import Compel, ReturnedEmbeddingsType
from diffusers import DiffusionPipeline
from diffusers.utils import make_image_grid
import torch
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
variant="fp16",
use_safetensors=True,
torch_dtype=torch.float16
).to("cuda")
compel = Compel(
tokenizer=[pipeline.tokenizer, pipeline.tokenizer_2] ,
text_encoder=[pipeline.text_encoder, pipeline.text_encoder_2],
returned_embeddings_type=ReturnedEmbeddingsType.PENULTIMATE_HIDDEN_STATES_NON_NORMALIZED,
requires_pooled=[False, True]
)
```
This time, let's upweight "ball" by a factor of 1.5 for the first prompt, and downweight "ball" by 0.6 for the second prompt. The [`StableDiffusionXLPipeline`] also requires [`pooled_prompt_embeds`](https://huggingface.co/docs/diffusers/en/api/pipelines/stable_diffusion/stable_diffusion_xl#diffusers.StableDiffusionXLInpaintPipeline.__call__.pooled_prompt_embeds) (and optionally [`negative_pooled_prompt_embeds`](https://huggingface.co/docs/diffusers/en/api/pipelines/stable_diffusion/stable_diffusion_xl#diffusers.StableDiffusionXLInpaintPipeline.__call__.negative_pooled_prompt_embeds)) so you should pass those to the pipeline along with the conditioning tensors:
```py
# apply weights
prompt = ["a red cat playing with a (ball)1.5", "a red cat playing with a (ball)0.6"]
conditioning, pooled = compel(prompt)
# generate image
generator = [torch.Generator().manual_seed(33) for _ in range(len(prompt))]
images = pipeline(prompt_embeds=conditioning, pooled_prompt_embeds=pooled, generator=generator, num_inference_steps=30).images
make_image_grid(images, rows=1, cols=2)
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/compel/sdxl_ball1.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">"a red cat playing with a (ball)1.5"</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/compel/sdxl_ball2.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">"a red cat playing with a (ball)0.6"</figcaption>
</div>
</div>

View File

@@ -106,7 +106,7 @@ Let's try it out!
## Deconstruct the Stable Diffusion pipeline
Stable Diffusion is a text-to-image *latent diffusion* model. It is called a latent diffusion model because it works with a lower-dimensional representation of the image instead of the actual pixel space, which makes it more memory efficient. The encoder compresses the image into a smaller representation, and a decoder converts the compressed representation back into an image. For text-to-image models, you'll need a tokenizer and an encoder to generate text embeddings. From the previous example, you already know you need a UNet model and a scheduler.
Stable Diffusion is a text-to-image *latent diffusion* model. It is called a latent diffusion model because it works with a lower-dimensional representation of the image instead of the actual pixel space, which makes it more memory efficient. The encoder compresses the image into a smaller representation, and a decoder to convert the compressed representation back into an image. For text-to-image models, you'll need a tokenizer and an encoder to generate text embeddings. From the previous example, you already know you need a UNet model and a scheduler.
As you can see, this is already more complex than the DDPM pipeline which only contains a UNet model. The Stable Diffusion model has three separate pretrained models.

View File

@@ -40,9 +40,9 @@ Training examples show how to pretrain or fine-tune diffusion models for a varie
| [**Text-to-Image fine-tuning**](./text_to_image) | ✅ | ✅ |
| [**Textual Inversion**](./textual_inversion) | ✅ | - | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/sd_textual_inversion_training.ipynb)
| [**Dreambooth**](./dreambooth) | ✅ | - | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/sd_dreambooth_training.ipynb)
| [**ControlNet**](./controlnet) | ✅ | ✅ | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/controlnet.ipynb)
| [**InstructPix2Pix**](./instruct_pix2pix) | ✅ | ✅ | [Notebook](https://github.com/huggingface/notebooks/blob/main/diffusers/InstructPix2Pix_using_diffusers.ipynb)
| [**Reinforcement Learning for Control**](./reinforcement_learning) | - | - | [Notebook1](https://github.com/huggingface/notebooks/blob/main/diffusers/reinforcement_learning_for_control.ipynb), [Notebook2](https://github.com/huggingface/notebooks/blob/main/diffusers/reinforcement_learning_with_diffusers.ipynb)
| [**ControlNet**](./controlnet) | ✅ | ✅ | -
| [**InstructPix2Pix**](./instruct_pix2pix) | ✅ | ✅ | -
| [**Reinforcement Learning for Control**](./reinforcement_learning) | - | - | coming soon.
## Community

View File

@@ -227,7 +227,7 @@ def log_validation(
pipeline.set_progress_bar_config(disable=True)
# run inference
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed is not None else None
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed else None
autocast_ctx = nullcontext()
with autocast_ctx:
@@ -880,7 +880,9 @@ class TokenEmbeddingsHandler:
idx_to_text_encoder_name = {0: "clip_l", 1: "t5"}
for idx, text_encoder in enumerate(self.text_encoders):
train_ids = self.train_ids if idx == 0 else self.train_ids_t5
embeds = text_encoder.text_model.embeddings.token_embedding if idx == 0 else text_encoder.shared
embeds = (
text_encoder.text_model.embeddings.token_embedding if idx == 0 else text_encoder.encoder.embed_tokens
)
assert embeds.weight.data.shape[0] == len(self.tokenizers[idx]), "Tokenizers should be the same."
new_token_embeddings = embeds.weight.data[train_ids]
@@ -902,7 +904,9 @@ class TokenEmbeddingsHandler:
@torch.no_grad()
def retract_embeddings(self):
for idx, text_encoder in enumerate(self.text_encoders):
embeds = text_encoder.text_model.embeddings.token_embedding if idx == 0 else text_encoder.shared
embeds = (
text_encoder.text_model.embeddings.token_embedding if idx == 0 else text_encoder.encoder.embed_tokens
)
index_no_updates = self.embeddings_settings[f"index_no_updates_{idx}"]
embeds.weight.data[index_no_updates] = (
self.embeddings_settings[f"original_embeddings_{idx}"][index_no_updates]
@@ -1745,7 +1749,7 @@ def main(args):
if args.enable_t5_ti: # whether to do pivotal tuning/textual inversion for T5 as well
text_lora_parameters_two = []
for name, param in text_encoder_two.named_parameters():
if "shared" in name:
if "token_embedding" in name:
# ensure that dtype is float32, even if rest of the model that isn't trained is loaded in fp16
param.data = param.to(dtype=torch.float32)
param.requires_grad = True

View File

@@ -1883,11 +1883,7 @@ def main(args):
pipeline.set_progress_bar_config(disable=True)
# run inference
generator = (
torch.Generator(device=accelerator.device).manual_seed(args.seed)
if args.seed is not None
else None
)
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed else None
pipeline_args = {"prompt": args.validation_prompt}
if torch.backends.mps.is_available():
@@ -1991,9 +1987,7 @@ def main(args):
)
# run inference
pipeline = pipeline.to(accelerator.device)
generator = (
torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed is not None else None
)
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed else None
images = [
pipeline(args.validation_prompt, num_inference_steps=25, generator=generator).images[0]
for _ in range(args.num_validation_images)

View File

@@ -269,7 +269,7 @@ def log_validation(
pipeline.set_progress_bar_config(disable=True)
# run inference
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed is not None else None
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed else None
# Currently the context determination is a bit hand-wavy. We can improve it in the future if there's a better
# way to condition it. Reference: https://github.com/huggingface/diffusers/pull/7126#issuecomment-1968523051
if torch.backends.mps.is_available() or "playground" in args.pretrained_model_name_or_path:

View File

@@ -722,7 +722,7 @@ def log_validation(
# pipe.set_progress_bar_config(disable=True)
# run inference
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed is not None else None
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed else None
videos = []
for _ in range(args.num_validation_videos):

View File

@@ -739,7 +739,7 @@ def log_validation(
# pipe.set_progress_bar_config(disable=True)
# run inference
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed is not None else None
generator = torch.Generator(device=accelerator.device).manual_seed(args.seed) if args.seed else None
videos = []
for _ in range(args.num_validation_videos):

File diff suppressed because it is too large Load Diff

View File

@@ -92,13 +92,9 @@ class CheckpointMergerPipeline(DiffusionPipeline):
token = kwargs.pop("token", None)
variant = kwargs.pop("variant", None)
revision = kwargs.pop("revision", None)
torch_dtype = kwargs.pop("torch_dtype", torch.float32)
torch_dtype = kwargs.pop("torch_dtype", None)
device_map = kwargs.pop("device_map", None)
if not isinstance(torch_dtype, torch.dtype):
torch_dtype = torch.float32
print(f"Passed `torch_dtype` {torch_dtype} is not a `torch.dtype`. Defaulting to `torch.float32`.")
alpha = kwargs.pop("alpha", 0.5)
interp = kwargs.pop("interp", None)

View File

@@ -1,645 +0,0 @@
"""
This script performs DDIM inversion for video frames using a pre-trained model and generates
a video reconstruction based on a provided prompt. It utilizes the CogVideoX pipeline to
process video frames, apply the DDIM inverse scheduler, and produce an output video.
**Please notice that this script is based on the CogVideoX 5B model, and would not generate
a good result for 2B variants.**
Usage:
python cogvideox_ddim_inversion.py
--model-path /path/to/model
--prompt "a prompt"
--video-path /path/to/video.mp4
--output-path /path/to/output
For more details about the cli arguments, please run `python cogvideox_ddim_inversion.py --help`.
Author:
LittleNyima <littlenyima[at]163[dot]com>
"""
import argparse
import math
import os
from typing import Any, Dict, List, Optional, Tuple, TypedDict, Union, cast
import torch
import torch.nn.functional as F
import torchvision.transforms as T
from transformers import T5EncoderModel, T5Tokenizer
from diffusers.models.attention_processor import Attention, CogVideoXAttnProcessor2_0
from diffusers.models.autoencoders import AutoencoderKLCogVideoX
from diffusers.models.embeddings import apply_rotary_emb
from diffusers.models.transformers.cogvideox_transformer_3d import CogVideoXBlock, CogVideoXTransformer3DModel
from diffusers.pipelines.cogvideo.pipeline_cogvideox import CogVideoXPipeline, retrieve_timesteps
from diffusers.schedulers import CogVideoXDDIMScheduler, DDIMInverseScheduler
from diffusers.utils import export_to_video
# Must import after torch because this can sometimes lead to a nasty segmentation fault, or stack smashing error.
# Very few bug reports but it happens. Look in decord Github issues for more relevant information.
import decord # isort: skip
class DDIMInversionArguments(TypedDict):
model_path: str
prompt: str
video_path: str
output_path: str
guidance_scale: float
num_inference_steps: int
skip_frames_start: int
skip_frames_end: int
frame_sample_step: Optional[int]
max_num_frames: int
width: int
height: int
fps: int
dtype: torch.dtype
seed: int
device: torch.device
def get_args() -> DDIMInversionArguments:
parser = argparse.ArgumentParser()
parser.add_argument("--model_path", type=str, required=True, help="Path of the pretrained model")
parser.add_argument("--prompt", type=str, required=True, help="Prompt for the direct sample procedure")
parser.add_argument("--video_path", type=str, required=True, help="Path of the video for inversion")
parser.add_argument("--output_path", type=str, default="output", help="Path of the output videos")
parser.add_argument("--guidance_scale", type=float, default=6.0, help="Classifier-free guidance scale")
parser.add_argument("--num_inference_steps", type=int, default=50, help="Number of inference steps")
parser.add_argument("--skip_frames_start", type=int, default=0, help="Number of skipped frames from the start")
parser.add_argument("--skip_frames_end", type=int, default=0, help="Number of skipped frames from the end")
parser.add_argument("--frame_sample_step", type=int, default=None, help="Temporal stride of the sampled frames")
parser.add_argument("--max_num_frames", type=int, default=81, help="Max number of sampled frames")
parser.add_argument("--width", type=int, default=720, help="Resized width of the video frames")
parser.add_argument("--height", type=int, default=480, help="Resized height of the video frames")
parser.add_argument("--fps", type=int, default=8, help="Frame rate of the output videos")
parser.add_argument("--dtype", type=str, default="bf16", choices=["bf16", "fp16"], help="Dtype of the model")
parser.add_argument("--seed", type=int, default=42, help="Seed for the random number generator")
parser.add_argument("--device", type=str, default="cuda", choices=["cuda", "cpu"], help="Device for inference")
args = parser.parse_args()
args.dtype = torch.bfloat16 if args.dtype == "bf16" else torch.float16
args.device = torch.device(args.device)
return DDIMInversionArguments(**vars(args))
class CogVideoXAttnProcessor2_0ForDDIMInversion(CogVideoXAttnProcessor2_0):
def __init__(self):
super().__init__()
def calculate_attention(
self,
query: torch.Tensor,
key: torch.Tensor,
value: torch.Tensor,
attn: Attention,
batch_size: int,
image_seq_length: int,
text_seq_length: int,
attention_mask: Optional[torch.Tensor],
image_rotary_emb: Optional[torch.Tensor],
) -> Tuple[torch.Tensor, torch.Tensor]:
r"""
Core attention computation with inversion-guided RoPE integration.
Args:
query (`torch.Tensor`): `[batch_size, seq_len, dim]` query tensor
key (`torch.Tensor`): `[batch_size, seq_len, dim]` key tensor
value (`torch.Tensor`): `[batch_size, seq_len, dim]` value tensor
attn (`Attention`): Parent attention module with projection layers
batch_size (`int`): Effective batch size (after chunk splitting)
image_seq_length (`int`): Length of image feature sequence
text_seq_length (`int`): Length of text feature sequence
attention_mask (`Optional[torch.Tensor]`): Attention mask tensor
image_rotary_emb (`Optional[torch.Tensor]`): Rotary embeddings for image positions
Returns:
`Tuple[torch.Tensor, torch.Tensor]`:
(1) hidden_states: [batch_size, image_seq_length, dim] processed image features
(2) encoder_hidden_states: [batch_size, text_seq_length, dim] processed text features
"""
inner_dim = key.shape[-1]
head_dim = inner_dim // attn.heads
query = query.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
key = key.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
value = value.view(batch_size, -1, attn.heads, head_dim).transpose(1, 2)
if attn.norm_q is not None:
query = attn.norm_q(query)
if attn.norm_k is not None:
key = attn.norm_k(key)
# Apply RoPE if needed
if image_rotary_emb is not None:
query[:, :, text_seq_length:] = apply_rotary_emb(query[:, :, text_seq_length:], image_rotary_emb)
if not attn.is_cross_attention:
if key.size(2) == query.size(2): # Attention for reference hidden states
key[:, :, text_seq_length:] = apply_rotary_emb(key[:, :, text_seq_length:], image_rotary_emb)
else: # RoPE should be applied to each group of image tokens
key[:, :, text_seq_length : text_seq_length + image_seq_length] = apply_rotary_emb(
key[:, :, text_seq_length : text_seq_length + image_seq_length], image_rotary_emb
)
key[:, :, text_seq_length * 2 + image_seq_length :] = apply_rotary_emb(
key[:, :, text_seq_length * 2 + image_seq_length :], image_rotary_emb
)
hidden_states = F.scaled_dot_product_attention(
query, key, value, attn_mask=attention_mask, dropout_p=0.0, is_causal=False
)
hidden_states = hidden_states.transpose(1, 2).reshape(batch_size, -1, attn.heads * head_dim)
# linear proj
hidden_states = attn.to_out[0](hidden_states)
# dropout
hidden_states = attn.to_out[1](hidden_states)
encoder_hidden_states, hidden_states = hidden_states.split(
[text_seq_length, hidden_states.size(1) - text_seq_length], dim=1
)
return hidden_states, encoder_hidden_states
def __call__(
self,
attn: Attention,
hidden_states: torch.Tensor,
encoder_hidden_states: torch.Tensor,
attention_mask: Optional[torch.Tensor] = None,
image_rotary_emb: Optional[torch.Tensor] = None,
) -> Tuple[torch.Tensor, torch.Tensor]:
r"""
Process the dual-path attention for the inversion-guided denoising procedure.
Args:
attn (`Attention`): Parent attention module
hidden_states (`torch.Tensor`): `[batch_size, image_seq_len, dim]` Image tokens
encoder_hidden_states (`torch.Tensor`): `[batch_size, text_seq_len, dim]` Text tokens
attention_mask (`Optional[torch.Tensor]`): Optional attention mask
image_rotary_emb (`Optional[torch.Tensor]`): Rotary embeddings for image tokens
Returns:
`Tuple[torch.Tensor, torch.Tensor]`:
(1) Final hidden states: `[batch_size, image_seq_length, dim]` Resulting image tokens
(2) Final encoder states: `[batch_size, text_seq_length, dim]` Resulting text tokens
"""
image_seq_length = hidden_states.size(1)
text_seq_length = encoder_hidden_states.size(1)
hidden_states = torch.cat([encoder_hidden_states, hidden_states], dim=1)
batch_size, sequence_length, _ = (
hidden_states.shape if encoder_hidden_states is None else encoder_hidden_states.shape
)
if attention_mask is not None:
attention_mask = attn.prepare_attention_mask(attention_mask, sequence_length, batch_size)
attention_mask = attention_mask.view(batch_size, attn.heads, -1, attention_mask.shape[-1])
query = attn.to_q(hidden_states)
key = attn.to_k(hidden_states)
value = attn.to_v(hidden_states)
query, query_reference = query.chunk(2)
key, key_reference = key.chunk(2)
value, value_reference = value.chunk(2)
batch_size = batch_size // 2
hidden_states, encoder_hidden_states = self.calculate_attention(
query=query,
key=torch.cat((key, key_reference), dim=1),
value=torch.cat((value, value_reference), dim=1),
attn=attn,
batch_size=batch_size,
image_seq_length=image_seq_length,
text_seq_length=text_seq_length,
attention_mask=attention_mask,
image_rotary_emb=image_rotary_emb,
)
hidden_states_reference, encoder_hidden_states_reference = self.calculate_attention(
query=query_reference,
key=key_reference,
value=value_reference,
attn=attn,
batch_size=batch_size,
image_seq_length=image_seq_length,
text_seq_length=text_seq_length,
attention_mask=attention_mask,
image_rotary_emb=image_rotary_emb,
)
return (
torch.cat((hidden_states, hidden_states_reference)),
torch.cat((encoder_hidden_states, encoder_hidden_states_reference)),
)
class OverrideAttnProcessors:
r"""
Context manager for temporarily overriding attention processors in CogVideo transformer blocks.
Designed for DDIM inversion process, replaces original attention processors with
`CogVideoXAttnProcessor2_0ForDDIMInversion` and restores them upon exit. Uses Python context manager
pattern to safely manage processor replacement.
Typical usage:
```python
with OverrideAttnProcessors(transformer):
# Perform DDIM inversion operations
```
Args:
transformer (`CogVideoXTransformer3DModel`):
The transformer model containing attention blocks to be modified. Should have
`transformer_blocks` attribute containing `CogVideoXBlock` instances.
"""
def __init__(self, transformer: CogVideoXTransformer3DModel):
self.transformer = transformer
self.original_processors = {}
def __enter__(self):
for block in self.transformer.transformer_blocks:
block = cast(CogVideoXBlock, block)
self.original_processors[id(block)] = block.attn1.get_processor()
block.attn1.set_processor(CogVideoXAttnProcessor2_0ForDDIMInversion())
def __exit__(self, _0, _1, _2):
for block in self.transformer.transformer_blocks:
block = cast(CogVideoXBlock, block)
block.attn1.set_processor(self.original_processors[id(block)])
def get_video_frames(
video_path: str,
width: int,
height: int,
skip_frames_start: int,
skip_frames_end: int,
max_num_frames: int,
frame_sample_step: Optional[int],
) -> torch.FloatTensor:
"""
Extract and preprocess video frames from a video file for VAE processing.
Args:
video_path (`str`): Path to input video file
width (`int`): Target frame width for decoding
height (`int`): Target frame height for decoding
skip_frames_start (`int`): Number of frames to skip at video start
skip_frames_end (`int`): Number of frames to skip at video end
max_num_frames (`int`): Maximum allowed number of output frames
frame_sample_step (`Optional[int]`):
Frame sampling step size. If None, automatically calculated as:
(total_frames - skipped_frames) // max_num_frames
Returns:
`torch.FloatTensor`: Preprocessed frames in `[F, C, H, W]` format where:
- `F`: Number of frames (adjusted to 4k + 1 for VAE compatibility)
- `C`: Channels (3 for RGB)
- `H`: Frame height
- `W`: Frame width
"""
with decord.bridge.use_torch():
video_reader = decord.VideoReader(uri=video_path, width=width, height=height)
video_num_frames = len(video_reader)
start_frame = min(skip_frames_start, video_num_frames)
end_frame = max(0, video_num_frames - skip_frames_end)
if end_frame <= start_frame:
indices = [start_frame]
elif end_frame - start_frame <= max_num_frames:
indices = list(range(start_frame, end_frame))
else:
step = frame_sample_step or (end_frame - start_frame) // max_num_frames
indices = list(range(start_frame, end_frame, step))
frames = video_reader.get_batch(indices=indices)
frames = frames[:max_num_frames].float() # ensure that we don't go over the limit
# Choose first (4k + 1) frames as this is how many is required by the VAE
selected_num_frames = frames.size(0)
remainder = (3 + selected_num_frames) % 4
if remainder != 0:
frames = frames[:-remainder]
assert frames.size(0) % 4 == 1
# Normalize the frames
transform = T.Lambda(lambda x: x / 255.0 * 2.0 - 1.0)
frames = torch.stack(tuple(map(transform, frames)), dim=0)
return frames.permute(0, 3, 1, 2).contiguous() # [F, C, H, W]
class CogVideoXDDIMInversionOutput:
inverse_latents: torch.FloatTensor
recon_latents: torch.FloatTensor
def __init__(self, inverse_latents: torch.FloatTensor, recon_latents: torch.FloatTensor):
self.inverse_latents = inverse_latents
self.recon_latents = recon_latents
class CogVideoXPipelineForDDIMInversion(CogVideoXPipeline):
def __init__(
self,
tokenizer: T5Tokenizer,
text_encoder: T5EncoderModel,
vae: AutoencoderKLCogVideoX,
transformer: CogVideoXTransformer3DModel,
scheduler: CogVideoXDDIMScheduler,
):
super().__init__(
tokenizer=tokenizer,
text_encoder=text_encoder,
vae=vae,
transformer=transformer,
scheduler=scheduler,
)
self.inverse_scheduler = DDIMInverseScheduler(**scheduler.config)
def encode_video_frames(self, video_frames: torch.FloatTensor) -> torch.FloatTensor:
"""
Encode video frames into latent space using Variational Autoencoder.
Args:
video_frames (`torch.FloatTensor`):
Input frames tensor in `[F, C, H, W]` format from `get_video_frames()`
Returns:
`torch.FloatTensor`: Encoded latents in `[1, F, D, H_latent, W_latent]` format where:
- `F`: Number of frames (same as input)
- `D`: Latent channel dimension
- `H_latent`: Latent space height (H // 2^vae.downscale_factor)
- `W_latent`: Latent space width (W // 2^vae.downscale_factor)
"""
vae: AutoencoderKLCogVideoX = self.vae
video_frames = video_frames.to(device=vae.device, dtype=vae.dtype)
video_frames = video_frames.unsqueeze(0).permute(0, 2, 1, 3, 4) # [B, C, F, H, W]
latent_dist = vae.encode(x=video_frames).latent_dist.sample().transpose(1, 2)
return latent_dist * vae.config.scaling_factor
@torch.no_grad()
def export_latents_to_video(self, latents: torch.FloatTensor, video_path: str, fps: int):
r"""
Decode latent vectors into video and export as video file.
Args:
latents (`torch.FloatTensor`): Encoded latents in `[B, F, D, H_latent, W_latent]` format from
`encode_video_frames()`
video_path (`str`): Output path for video file
fps (`int`): Target frames per second for output video
"""
video = self.decode_latents(latents)
frames = self.video_processor.postprocess_video(video=video, output_type="pil")
os.makedirs(os.path.dirname(video_path), exist_ok=True)
export_to_video(video_frames=frames[0], output_video_path=video_path, fps=fps)
# Modified from CogVideoXPipeline.__call__
@torch.no_grad()
def sample(
self,
latents: torch.FloatTensor,
scheduler: Union[DDIMInverseScheduler, CogVideoXDDIMScheduler],
prompt: Optional[Union[str, List[str]]] = None,
negative_prompt: Optional[Union[str, List[str]]] = None,
num_inference_steps: int = 50,
guidance_scale: float = 6,
use_dynamic_cfg: bool = False,
eta: float = 0.0,
generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None,
attention_kwargs: Optional[Dict[str, Any]] = None,
reference_latents: torch.FloatTensor = None,
) -> torch.FloatTensor:
r"""
Execute the core sampling loop for video generation/inversion using CogVideoX.
Implements the full denoising trajectory recording for both DDIM inversion and
generation processes. Supports dynamic classifier-free guidance and reference
latent conditioning.
Args:
latents (`torch.FloatTensor`):
Initial noise tensor of shape `[B, F, C, H, W]`.
scheduler (`Union[DDIMInverseScheduler, CogVideoXDDIMScheduler]`):
Scheduling strategy for diffusion process. Use:
(1) `DDIMInverseScheduler` for inversion
(2) `CogVideoXDDIMScheduler` for generation
prompt (`Optional[Union[str, List[str]]]`):
Text prompt(s) for conditional generation. Defaults to unconditional.
negative_prompt (`Optional[Union[str, List[str]]]`):
Negative prompt(s) for guidance. Requires `guidance_scale > 1`.
num_inference_steps (`int`):
Number of denoising steps. Affects quality/compute trade-off.
guidance_scale (`float`):
Classifier-free guidance weight. 1.0 = no guidance.
use_dynamic_cfg (`bool`):
Enable time-varying guidance scale (cosine schedule)
eta (`float`):
DDIM variance parameter (0 = deterministic process)
generator (`Optional[Union[torch.Generator, List[torch.Generator]]]`):
Random number generator(s) for reproducibility
attention_kwargs (`Optional[Dict[str, Any]]`):
Custom parameters for attention modules
reference_latents (`torch.FloatTensor`):
Reference latent trajectory for conditional sampling. Shape should match
`[T, B, F, C, H, W]` where `T` is number of timesteps
Returns:
`torch.FloatTensor`:
Full denoising trajectory tensor of shape `[T, B, F, C, H, W]`.
"""
self._guidance_scale = guidance_scale
self._attention_kwargs = attention_kwargs
self._interrupt = False
device = self._execution_device
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
# corresponds to doing no classifier free guidance.
do_classifier_free_guidance = guidance_scale > 1.0
# 3. Encode input prompt
prompt_embeds, negative_prompt_embeds = self.encode_prompt(
prompt,
negative_prompt,
do_classifier_free_guidance,
device=device,
)
if do_classifier_free_guidance:
prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds], dim=0)
if reference_latents is not None:
prompt_embeds = torch.cat([prompt_embeds] * 2, dim=0)
# 4. Prepare timesteps
timesteps, num_inference_steps = retrieve_timesteps(scheduler, num_inference_steps, device)
self._num_timesteps = len(timesteps)
# 5. Prepare latents.
latents = latents.to(device=device) * scheduler.init_noise_sigma
# 6. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
if isinstance(scheduler, DDIMInverseScheduler): # Inverse scheduler does not accept extra kwargs
extra_step_kwargs = {}
# 7. Create rotary embeds if required
image_rotary_emb = (
self._prepare_rotary_positional_embeddings(
height=latents.size(3) * self.vae_scale_factor_spatial,
width=latents.size(4) * self.vae_scale_factor_spatial,
num_frames=latents.size(1),
device=device,
)
if self.transformer.config.use_rotary_positional_embeddings
else None
)
# 8. Denoising loop
num_warmup_steps = max(len(timesteps) - num_inference_steps * scheduler.order, 0)
trajectory = torch.zeros_like(latents).unsqueeze(0).repeat(len(timesteps), 1, 1, 1, 1, 1)
with self.progress_bar(total=num_inference_steps) as progress_bar:
for i, t in enumerate(timesteps):
if self.interrupt:
continue
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
if reference_latents is not None:
reference = reference_latents[i]
reference = torch.cat([reference] * 2) if do_classifier_free_guidance else reference
latent_model_input = torch.cat([latent_model_input, reference], dim=0)
latent_model_input = scheduler.scale_model_input(latent_model_input, t)
# broadcast to batch dimension in a way that's compatible with ONNX/Core ML
timestep = t.expand(latent_model_input.shape[0])
# predict noise model_output
noise_pred = self.transformer(
hidden_states=latent_model_input,
encoder_hidden_states=prompt_embeds,
timestep=timestep,
image_rotary_emb=image_rotary_emb,
attention_kwargs=attention_kwargs,
return_dict=False,
)[0]
noise_pred = noise_pred.float()
if reference_latents is not None: # Recover the original batch size
noise_pred, _ = noise_pred.chunk(2)
# perform guidance
if use_dynamic_cfg:
self._guidance_scale = 1 + guidance_scale * (
(1 - math.cos(math.pi * ((num_inference_steps - t.item()) / num_inference_steps) ** 5.0)) / 2
)
if do_classifier_free_guidance:
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + self.guidance_scale * (noise_pred_text - noise_pred_uncond)
# compute the noisy sample x_t-1 -> x_t
latents = scheduler.step(noise_pred, t, latents, **extra_step_kwargs, return_dict=False)[0]
latents = latents.to(prompt_embeds.dtype)
trajectory[i] = latents
if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % scheduler.order == 0):
progress_bar.update()
# Offload all models
self.maybe_free_model_hooks()
return trajectory
@torch.no_grad()
def __call__(
self,
prompt: str,
video_path: str,
guidance_scale: float,
num_inference_steps: int,
skip_frames_start: int,
skip_frames_end: int,
frame_sample_step: Optional[int],
max_num_frames: int,
width: int,
height: int,
seed: int,
):
"""
Performs DDIM inversion on a video to reconstruct it with a new prompt.
Args:
prompt (`str`): The text prompt to guide the reconstruction.
video_path (`str`): Path to the input video file.
guidance_scale (`float`): Scale for classifier-free guidance.
num_inference_steps (`int`): Number of denoising steps.
skip_frames_start (`int`): Number of frames to skip from the beginning of the video.
skip_frames_end (`int`): Number of frames to skip from the end of the video.
frame_sample_step (`Optional[int]`): Step size for sampling frames. If None, all frames are used.
max_num_frames (`int`): Maximum number of frames to process.
width (`int`): Width of the output video frames.
height (`int`): Height of the output video frames.
seed (`int`): Random seed for reproducibility.
Returns:
`CogVideoXDDIMInversionOutput`: Contains the inverse latents and reconstructed latents.
"""
if not self.transformer.config.use_rotary_positional_embeddings:
raise NotImplementedError("This script supports CogVideoX 5B model only.")
video_frames = get_video_frames(
video_path=video_path,
width=width,
height=height,
skip_frames_start=skip_frames_start,
skip_frames_end=skip_frames_end,
max_num_frames=max_num_frames,
frame_sample_step=frame_sample_step,
).to(device=self.device)
video_latents = self.encode_video_frames(video_frames=video_frames)
inverse_latents = self.sample(
latents=video_latents,
scheduler=self.inverse_scheduler,
prompt="",
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
generator=torch.Generator(device=self.device).manual_seed(seed),
)
with OverrideAttnProcessors(transformer=self.transformer):
recon_latents = self.sample(
latents=torch.randn_like(video_latents),
scheduler=self.scheduler,
prompt=prompt,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
generator=torch.Generator(device=self.device).manual_seed(seed),
reference_latents=reversed(inverse_latents),
)
return CogVideoXDDIMInversionOutput(
inverse_latents=inverse_latents,
recon_latents=recon_latents,
)
if __name__ == "__main__":
arguments = get_args()
pipeline = CogVideoXPipelineForDDIMInversion.from_pretrained(
arguments.pop("model_path"),
torch_dtype=arguments.pop("dtype"),
).to(device=arguments.pop("device"))
output_path = arguments.pop("output_path")
fps = arguments.pop("fps")
inverse_video_path = os.path.join(output_path, f"{arguments.get('video_path')}_inversion.mp4")
recon_video_path = os.path.join(output_path, f"{arguments.get('video_path')}_reconstruction.mp4")
# Run DDIM inversion
output = pipeline(**arguments)
pipeline.export_latents_to_video(output.inverse_latents[-1], inverse_video_path, fps)
pipeline.export_latents_to_video(output.recon_latents[-1], recon_video_path, fps)

File diff suppressed because it is too large Load Diff

File diff suppressed because it is too large Load Diff

View File

@@ -87,7 +87,7 @@ def calculate_shift(
base_seq_len: int = 256,
max_seq_len: int = 4096,
base_shift: float = 0.5,
max_shift: float = 1.15,
max_shift: float = 1.16,
):
m = (max_shift - base_shift) / (max_seq_len - base_seq_len)
b = base_shift - m * base_seq_len
@@ -878,7 +878,7 @@ class FluxDifferentialImg2ImgPipeline(DiffusionPipeline, FluxLoraLoaderMixin):
self.scheduler.config.get("base_image_seq_len", 256),
self.scheduler.config.get("max_image_seq_len", 4096),
self.scheduler.config.get("base_shift", 0.5),
self.scheduler.config.get("max_shift", 1.15),
self.scheduler.config.get("max_shift", 1.16),
)
timesteps, num_inference_steps = retrieve_timesteps(
self.scheduler,

View File

@@ -94,7 +94,7 @@ def calculate_shift(
base_seq_len: int = 256,
max_seq_len: int = 4096,
base_shift: float = 0.5,
max_shift: float = 1.15,
max_shift: float = 1.16,
):
m = (max_shift - base_shift) / (max_seq_len - base_seq_len)
b = base_shift - m * base_seq_len
@@ -823,7 +823,7 @@ class RFInversionFluxPipeline(
self.scheduler.config.get("base_image_seq_len", 256),
self.scheduler.config.get("max_image_seq_len", 4096),
self.scheduler.config.get("base_shift", 0.5),
self.scheduler.config.get("max_shift", 1.15),
self.scheduler.config.get("max_shift", 1.16),
)
timesteps, num_inference_steps = retrieve_timesteps(
self.scheduler,
@@ -993,7 +993,7 @@ class RFInversionFluxPipeline(
self.scheduler.config.get("base_image_seq_len", 256),
self.scheduler.config.get("max_image_seq_len", 4096),
self.scheduler.config.get("base_shift", 0.5),
self.scheduler.config.get("max_shift", 1.15),
self.scheduler.config.get("max_shift", 1.16),
)
timesteps, num_inversion_steps = retrieve_timesteps(
self.scheduler,

View File

@@ -91,7 +91,7 @@ def calculate_shift(
base_seq_len: int = 256,
max_seq_len: int = 4096,
base_shift: float = 0.5,
max_shift: float = 1.15,
max_shift: float = 1.16,
):
m = (max_shift - base_shift) / (max_seq_len - base_seq_len)
b = base_shift - m * base_seq_len
@@ -1041,7 +1041,7 @@ class FluxSemanticGuidancePipeline(
self.scheduler.config.get("base_image_seq_len", 256),
self.scheduler.config.get("max_image_seq_len", 4096),
self.scheduler.config.get("base_shift", 0.5),
self.scheduler.config.get("max_shift", 1.15),
self.scheduler.config.get("max_shift", 1.16),
)
timesteps, num_inference_steps = retrieve_timesteps(
self.scheduler,

View File

@@ -70,7 +70,7 @@ def calculate_shift(
base_seq_len: int = 256,
max_seq_len: int = 4096,
base_shift: float = 0.5,
max_shift: float = 1.15,
max_shift: float = 1.16,
):
m = (max_shift - base_shift) / (max_seq_len - base_seq_len)
b = base_shift - m * base_seq_len
@@ -759,7 +759,7 @@ class FluxCFGPipeline(DiffusionPipeline, FluxLoraLoaderMixin, FromSingleFileMixi
self.scheduler.config.get("base_image_seq_len", 256),
self.scheduler.config.get("max_image_seq_len", 4096),
self.scheduler.config.get("base_shift", 0.5),
self.scheduler.config.get("max_shift", 1.15),
self.scheduler.config.get("max_shift", 1.16),
)
timesteps, num_inference_steps = retrieve_timesteps(
self.scheduler,

Some files were not shown because too many files have changed in this diff Show More