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1c693f9b68 |
@@ -280,7 +280,7 @@ init_image = init_image.resize((768, 512))
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prompt = "A fantasy landscape, trending on artstation"
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images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
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images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
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images[0].save("fantasy_landscape.png")
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```
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@@ -49,6 +49,8 @@
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title: "OpenVINO"
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- local: optimization/mps
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title: "MPS"
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- local: optimization/habana
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title: "Habana Gaudi"
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title: "Optimization/Special Hardware"
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- sections:
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- local: training/overview
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@@ -106,6 +108,8 @@
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title: "Score SDE VE"
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- local: api/pipelines/stable_diffusion
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title: "Stable Diffusion"
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- local: api/pipelines/stable_diffusion_2
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title: "Stable Diffusion 2"
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- local: api/pipelines/stable_diffusion_safe
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title: "Safe Stable Diffusion"
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- local: api/pipelines/stochastic_karras_ve
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@@ -51,7 +51,7 @@ To use a different scheduler, you can either change it via the [`ConfigMixin.fro
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```
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||||
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||||
|
||||
- *How to conver all use cases with multiple or single pipeline*
|
||||
- *How to convert all use cases with multiple or single pipeline*
|
||||
|
||||
If you want to use all possible use cases in a single `DiffusionPipeline` we recommend using the `components` functionality to instantiate all components in the most memory-efficient way:
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@@ -57,7 +57,7 @@ prompt = "An astronaut riding an elephant"
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image = pipe(
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prompt=prompt,
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source_prompt=source_prompt,
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init_image=init_image,
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image=init_image,
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num_inference_steps=100,
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eta=0.1,
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strength=0.8,
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@@ -83,7 +83,7 @@ torch.manual_seed(0)
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image = pipe(
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prompt=prompt,
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source_prompt=source_prompt,
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init_image=init_image,
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image=init_image,
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num_inference_steps=100,
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eta=0.1,
|
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strength=0.85,
|
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|
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@@ -58,6 +58,9 @@ available a colab notebook to directly try them out.
|
||||
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
|
||||
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
|
||||
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
|
||||
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-to-Image Generation |
|
||||
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting |
|
||||
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
|
||||
| [stable_diffusion_safe](./api/pipelines/stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
|
||||
| [stochastic_karras_ve](./api/pipelines/stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
|
||||
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
|
||||
@@ -146,7 +149,7 @@ init_image = init_image.resize((768, 512))
|
||||
|
||||
prompt = "A fantasy landscape, trending on artstation"
|
||||
|
||||
images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
|
||||
images[0].save("fantasy_landscape.png")
|
||||
```
|
||||
|
||||
@@ -48,7 +48,7 @@ To use a different scheduler, you can either change it via the [`ConfigMixin.fro
|
||||
```
|
||||
|
||||
|
||||
### How to conver all use cases with multiple or single pipeline
|
||||
### How to convert all use cases with multiple or single pipeline
|
||||
|
||||
If you want to use all possible use cases in a single `DiffusionPipeline` you can either:
|
||||
- Make use of the [Stable Diffusion Mega Pipeline](https://github.com/huggingface/diffusers/tree/main/examples/community#stable-diffusion-mega) or
|
||||
@@ -76,22 +76,40 @@ If you want to use all possible use cases in a single `DiffusionPipeline` you ca
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_vae_slicing
|
||||
- disable_vae_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
|
||||
## StableDiffusionImg2ImgPipeline
|
||||
[[autodoc]] StableDiffusionImg2ImgPipeline
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
|
||||
## StableDiffusionInpaintPipeline
|
||||
[[autodoc]] StableDiffusionInpaintPipeline
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
|
||||
## StableDiffusionImageVariationPipeline
|
||||
[[autodoc]] StableDiffusionImageVariationPipeline
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
|
||||
|
||||
## StableDiffusionUpscalePipeline
|
||||
[[autodoc]] StableDiffusionUpscalePipeline
|
||||
- __call__
|
||||
- enable_attention_slicing
|
||||
- disable_attention_slicing
|
||||
- enable_xformers_memory_efficient_attention
|
||||
- disable_xformers_memory_efficient_attention
|
||||
|
||||
142
docs/source/api/pipelines/stable_diffusion_2.mdx
Normal file
142
docs/source/api/pipelines/stable_diffusion_2.mdx
Normal file
@@ -0,0 +1,142 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# Stable diffusion 2
|
||||
|
||||
Stable Diffusion 2 is a text-to-image _latent diffusion_ model built upon the work of [Stable Diffusion 1](https://stability.ai/blog/stable-diffusion-public-release).
|
||||
The project to train Stable Diffusion 2 was led by Robin Rombach and Katherine Crowson from [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/).
|
||||
|
||||
*The Stable Diffusion 2.0 release includes robust text-to-image models trained using a brand new text encoder (OpenCLIP), developed by LAION with support from Stability AI, which greatly improves the quality of the generated images compared to earlier V1 releases. The text-to-image models in this release can generate images with default resolutions of both 512x512 pixels and 768x768 pixels.
|
||||
These models are trained on an aesthetic subset of the [LAION-5B dataset](https://laion.ai/blog/laion-5b/) created by the DeepFloyd team at Stability AI, which is then further filtered to remove adult content using [LAION’s NSFW filter](https://openreview.net/forum?id=M3Y74vmsMcY).*
|
||||
|
||||
For more details about how Stable Diffusion 2 works and how it differs from Stable Diffusion 1, please refer to the official [launch announcement post](https://stability.ai/blog/stable-diffusion-v2-release).
|
||||
|
||||
## Tips
|
||||
|
||||
### Available checkpoints:
|
||||
|
||||
Note that the architecture is more or less identical to [Stable Diffusion 1](./api/pipelines/stable_diffusion) so please refer to [this page](./api/pipelines/stable_diffusion) for API documentation.
|
||||
|
||||
- *Text-to-Image (512x512 resolution)*: [stabilityai/stable-diffusion-2-base](https://huggingface.co/stabilityai/stable-diffusion-2-base) with [`StableDiffusionPipeline`]
|
||||
- *Text-to-Image (768x768 resolution)*: [stabilityai/stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2) with [`StableDiffusionPipeline`]
|
||||
- *Image Inpainting (512x512 resolution)*: [stabilityai/stable-diffusion-2-inpainting](https://huggingface.co/stabilityai/stable-diffusion-2-inpainting) with [`StableDiffusionInpaintPipeline`]
|
||||
- *Image Upscaling (x4 resolution resolution)*: [stable-diffusion-x4-upscaler](https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler) [`StableDiffusionUpscalePipeline`]
|
||||
|
||||
We recommend using the [`DPMSolverMultistepScheduler`] as it's currently the fastest scheduler there is.
|
||||
|
||||
- *Text-to-Image (512x512 resolution)*:
|
||||
|
||||
```python
|
||||
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
|
||||
import torch
|
||||
|
||||
repo_id = "stabilityai/stable-diffusion-2-base"
|
||||
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
|
||||
|
||||
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "High quality photo of an astronaut riding a horse in space"
|
||||
image = pipe(prompt, num_inference_steps=25).images[0]
|
||||
image.save("astronaut.png")
|
||||
```
|
||||
|
||||
- *Text-to-Image (768x768 resolution)*:
|
||||
|
||||
```python
|
||||
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
|
||||
import torch
|
||||
|
||||
repo_id = "stabilityai/stable-diffusion-2"
|
||||
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
|
||||
|
||||
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "High quality photo of an astronaut riding a horse in space"
|
||||
image = pipe(prompt, guidance_scale=9, num_inference_steps=25).images[0]
|
||||
image.save("astronaut.png")
|
||||
```
|
||||
|
||||
- *Image Inpainting (512x512 resolution)*:
|
||||
|
||||
```python
|
||||
import PIL
|
||||
import requests
|
||||
import torch
|
||||
from io import BytesIO
|
||||
|
||||
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
|
||||
|
||||
|
||||
def download_image(url):
|
||||
response = requests.get(url)
|
||||
return PIL.Image.open(BytesIO(response.content)).convert("RGB")
|
||||
|
||||
|
||||
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
|
||||
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
|
||||
|
||||
init_image = download_image(img_url).resize((512, 512))
|
||||
mask_image = download_image(mask_url).resize((512, 512))
|
||||
|
||||
repo_id = "stabilityai/stable-diffusion-2-inpainting"
|
||||
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
|
||||
|
||||
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "Face of a yellow cat, high resolution, sitting on a park bench"
|
||||
image = pipe(prompt=prompt, image=init_image, mask_image=mask_image, num_inference_steps=25).images[0]
|
||||
|
||||
image.save("yellow_cat.png")
|
||||
```
|
||||
|
||||
- *Image Upscaling (x4 resolution resolution)*: [stable-diffusion-x4-upscaler](https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler) [`StableDiffusionUpscalePipeline`]
|
||||
|
||||
```python
|
||||
import requests
|
||||
from PIL import Image
|
||||
from io import BytesIO
|
||||
from diffusers import StableDiffusionUpscalePipeline
|
||||
import torch
|
||||
|
||||
# load model and scheduler
|
||||
model_id = "stabilityai/stable-diffusion-x4-upscaler"
|
||||
pipeline = StableDiffusionUpscalePipeline.from_pretrained(model_id, revision="fp16", torch_dtype=torch.float16)
|
||||
pipeline = pipeline.to("cuda")
|
||||
|
||||
# let's download an image
|
||||
url = "https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd2-upscale/low_res_cat.png"
|
||||
response = requests.get(url)
|
||||
low_res_img = Image.open(BytesIO(response.content)).convert("RGB")
|
||||
low_res_img = low_res_img.resize((128, 128))
|
||||
prompt = "a white cat"
|
||||
upscaled_image = pipeline(prompt=prompt, image=low_res_img).images[0]
|
||||
upscaled_image.save("upsampled_cat.png")
|
||||
```
|
||||
|
||||
### How to load and use different schedulers.
|
||||
|
||||
The stable diffusion pipeline uses [`DDIMScheduler`] scheduler by default. But `diffusers` provides many other schedulers that can be used with the stable diffusion pipeline such as [`PNDMScheduler`], [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`] etc.
|
||||
To use a different scheduler, you can either change it via the [`ConfigMixin.from_config`] method or pass the `scheduler` argument to the `from_pretrained` method of the pipeline. For example, to use the [`EulerDiscreteScheduler`], you can do the following:
|
||||
|
||||
```python
|
||||
>>> from diffusers import StableDiffusionPipeline, EulerDiscreteScheduler
|
||||
|
||||
>>> pipeline = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2")
|
||||
>>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
|
||||
|
||||
>>> # or
|
||||
>>> euler_scheduler = EulerDiscreteScheduler.from_pretrained("stabilityai/stable-diffusion-2", subfolder="scheduler")
|
||||
>>> pipeline = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2", scheduler=euler_scheduler)
|
||||
```
|
||||
@@ -48,6 +48,9 @@ available a colab notebook to directly try them out.
|
||||
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
|
||||
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
|
||||
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
|
||||
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-to-Image Generation |
|
||||
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting |
|
||||
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
|
||||
| [stable_diffusion_safe](./api/pipelines/stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
|
||||
| [stochastic_karras_ve](./api/pipelines/stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
|
||||
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
|
||||
|
||||
@@ -117,6 +117,34 @@ image = pipe(prompt).images[0]
|
||||
|
||||
There's a small performance penalty of about 10% slower inference times, but this method allows you to use Stable Diffusion in as little as 3.2 GB of VRAM!
|
||||
|
||||
|
||||
## Sliced VAE decode for larger batches
|
||||
|
||||
To decode large batches of images with limited VRAM, or to enable batches with 32 images or more, you can use sliced VAE decode that decodes the batch latents one image at a time.
|
||||
|
||||
You likely want to couple this with [`~StableDiffusionPipeline.enable_attention_slicing`] or [`~StableDiffusionPipeline.enable_xformers_memory_efficient_attention`] to further minimize memory use.
|
||||
|
||||
To perform the VAE decode one image at a time, invoke [`~StableDiffusionPipeline.enable_vae_slicing`] in your pipeline before inference. For example:
|
||||
|
||||
```Python
|
||||
import torch
|
||||
from diffusers import StableDiffusionPipeline
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained(
|
||||
"runwayml/stable-diffusion-v1-5",
|
||||
revision="fp16",
|
||||
torch_dtype=torch.float16,
|
||||
)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
prompt = "a photo of an astronaut riding a horse on mars"
|
||||
pipe.enable_vae_slicing()
|
||||
images = pipe([prompt] * 32).images
|
||||
```
|
||||
|
||||
You may see a small performance boost in VAE decode on multi-image batches. There should be no performance impact on single-image batches.
|
||||
|
||||
|
||||
## Offloading to CPU with accelerate for memory savings
|
||||
|
||||
For additional memory savings, you can offload the weights to CPU and load them to GPU when performing the forward pass.
|
||||
|
||||
70
docs/source/optimization/habana.mdx
Normal file
70
docs/source/optimization/habana.mdx
Normal file
@@ -0,0 +1,70 @@
|
||||
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
|
||||
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
|
||||
the License. You may obtain a copy of the License at
|
||||
|
||||
http://www.apache.org/licenses/LICENSE-2.0
|
||||
|
||||
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
|
||||
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
|
||||
specific language governing permissions and limitations under the License.
|
||||
-->
|
||||
|
||||
# How to use Stable Diffusion on Habana Gaudi
|
||||
|
||||
🤗 Diffusers is compatible with Habana Gaudi through 🤗 [Optimum Habana](https://huggingface.co/docs/optimum/habana/usage_guides/stable_diffusion).
|
||||
|
||||
## Requirements
|
||||
|
||||
- Optimum Habana 1.3 or later, [here](https://huggingface.co/docs/optimum/habana/installation) is how to install it.
|
||||
- SynapseAI 1.7.
|
||||
|
||||
|
||||
## Inference Pipeline
|
||||
|
||||
To generate images with Stable Diffusion 1 and 2 on Gaudi, you need to instantiate two instances:
|
||||
- A pipeline with [`GaudiStableDiffusionPipeline`](https://huggingface.co/docs/optimum/habana/package_reference/stable_diffusion_pipeline). This pipeline supports *text-to-image generation*.
|
||||
- A scheduler with [`GaudiDDIMScheduler`](https://huggingface.co/docs/optimum/habana/package_reference/stable_diffusion_pipeline#optimum.habana.diffusers.GaudiDDIMScheduler). This scheduler has been optimized for Habana Gaudi.
|
||||
|
||||
When initializing the pipeline, you have to specify `use_habana=True` to deploy it on HPUs.
|
||||
Furthermore, in order to get the fastest possible generations you should enable **HPU graphs** with `use_hpu_graphs=True`.
|
||||
Finally, you will need to specify a [Gaudi configuration](https://huggingface.co/docs/optimum/habana/package_reference/gaudi_config) which can be downloaded from the [Hugging Face Hub](https://huggingface.co/Habana).
|
||||
|
||||
```python
|
||||
from optimum.habana import GaudiConfig
|
||||
from optimum.habana.diffusers import GaudiDDIMScheduler, GaudiStableDiffusionPipeline
|
||||
|
||||
model_name = "stabilityai/stable-diffusion-2-base"
|
||||
scheduler = GaudiDDIMScheduler.from_pretrained(model_name, subfolder="scheduler")
|
||||
pipeline = GaudiStableDiffusionPipeline.from_pretrained(
|
||||
model_name,
|
||||
scheduler=scheduler,
|
||||
use_habana=True,
|
||||
use_hpu_graphs=True,
|
||||
gaudi_config="Habana/stable-diffusion",
|
||||
)
|
||||
```
|
||||
|
||||
You can then call the pipeline to generate images by batches from one or several prompts:
|
||||
```python
|
||||
outputs = pipeline(
|
||||
prompt=[
|
||||
"High quality photo of an astronaut riding a horse in space",
|
||||
"Face of a yellow cat, high resolution, sitting on a park bench",
|
||||
],
|
||||
num_images_per_prompt=10,
|
||||
batch_size=4,
|
||||
)
|
||||
```
|
||||
|
||||
For more information, check out Optimum Habana's [documentation](https://huggingface.co/docs/optimum/habana/usage_guides/stable_diffusion) and the [example](https://github.com/huggingface/optimum-habana/tree/main/examples/stable-diffusion) provided in the official Github repository.
|
||||
|
||||
|
||||
## Benchmark
|
||||
|
||||
Here are the latencies for Habana Gaudi 1 and Gaudi 2 with the [Habana/stable-diffusion](https://huggingface.co/Habana/stable-diffusion) Gaudi configuration (mixed precision bf16/fp32):
|
||||
|
||||
| | Latency | Batch size |
|
||||
| ------- |:-------:|:----------:|
|
||||
| Gaudi 1 | 4.37s | 4/8 |
|
||||
| Gaudi 2 | 1.19s | 4/8 |
|
||||
@@ -177,7 +177,7 @@ init_image = download_image(
|
||||
|
||||
prompt = "A fantasy landscape, trending on artstation"
|
||||
|
||||
images = pipe.img2img(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
images = pipe.img2img(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
|
||||
### Inpainting
|
||||
|
||||
@@ -187,7 +187,7 @@ init_image = download_image(img_url).resize((512, 512))
|
||||
mask_image = download_image(mask_url).resize((512, 512))
|
||||
|
||||
prompt = "a cat sitting on a bench"
|
||||
images = pipe.inpaint(prompt=prompt, init_image=init_image, mask_image=mask_image, strength=0.75).images
|
||||
images = pipe.inpaint(prompt=prompt, image=init_image, mask_image=mask_image, strength=0.75).images
|
||||
```
|
||||
|
||||
As shown above this one pipeline can run all both "text-to-image", "image-to-image", and "inpainting" in one pipeline.
|
||||
|
||||
@@ -37,7 +37,7 @@ init_image.thumbnail((768, 768))
|
||||
|
||||
prompt = "A fantasy landscape, trending on artstation"
|
||||
|
||||
images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
|
||||
images[0].save("fantasy_landscape.png")
|
||||
```
|
||||
|
||||
@@ -378,21 +378,3 @@ dpm = DPMSolverMultistepScheduler.from_pretrained(repo_id, subfolder="scheduler"
|
||||
# replace `dpm` with any of `ddpm`, `ddim`, `pndm`, `lms`, `euler`, `euler_anc`
|
||||
pipeline = StableDiffusionPipeline.from_pretrained(repo_id, scheduler=dpm)
|
||||
```
|
||||
|
||||
## API
|
||||
|
||||
[[autodoc]] modeling_utils.ModelMixin
|
||||
- from_pretrained
|
||||
- save_pretrained
|
||||
|
||||
[[autodoc]] pipeline_utils.DiffusionPipeline
|
||||
- from_pretrained
|
||||
- save_pretrained
|
||||
|
||||
[[autodoc]] modeling_flax_utils.FlaxModelMixin
|
||||
- from_pretrained
|
||||
- save_pretrained
|
||||
|
||||
[[autodoc]] pipeline_flax_utils.FlaxDiffusionPipeline
|
||||
- from_pretrained
|
||||
- save_pretrained
|
||||
|
||||
@@ -166,7 +166,7 @@ init_image = download_image("https://raw.githubusercontent.com/CompVis/stable-di
|
||||
|
||||
prompt = "A fantasy landscape, trending on artstation"
|
||||
|
||||
images = pipe.img2img(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
images = pipe.img2img(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
|
||||
### Inpainting
|
||||
|
||||
@@ -176,7 +176,7 @@ init_image = download_image(img_url).resize((512, 512))
|
||||
mask_image = download_image(mask_url).resize((512, 512))
|
||||
|
||||
prompt = "a cat sitting on a bench"
|
||||
images = pipe.inpaint(prompt=prompt, init_image=init_image, mask_image=mask_image, strength=0.75).images
|
||||
images = pipe.inpaint(prompt=prompt, image=init_image, mask_image=mask_image, strength=0.75).images
|
||||
```
|
||||
|
||||
As shown above this one pipeline can run all both "text-to-image", "image-to-image", and "inpainting" in one pipeline.
|
||||
@@ -420,7 +420,7 @@ init_image = Image.open(BytesIO(response.content)).convert("RGB")
|
||||
init_image = init_image.resize((512, 512))
|
||||
res = pipe.train(
|
||||
prompt,
|
||||
init_image,
|
||||
image=init_image,
|
||||
guidance_scale=7.5,
|
||||
num_inference_steps=50,
|
||||
generator=generator)
|
||||
@@ -602,7 +602,7 @@ For example, this could be used to place a logo on a shirt and make it blend sea
|
||||
import PIL
|
||||
import torch
|
||||
|
||||
from diffusers import StableDiffusionInpaintPipeline
|
||||
from diffusers import DiffusionPipeline
|
||||
|
||||
image_path = "./path-to-image.png"
|
||||
inner_image_path = "./path-to-inner-image.png"
|
||||
@@ -612,10 +612,11 @@ init_image = PIL.Image.open(image_path).convert("RGB").resize((512, 512))
|
||||
inner_image = PIL.Image.open(inner_image_path).convert("RGBA").resize((512, 512))
|
||||
mask_image = PIL.Image.open(mask_path).convert("RGB").resize((512, 512))
|
||||
|
||||
pipe = StableDiffusionInpaintPipeline.from_pretrained(
|
||||
pipe = DiffusionPipeline.from_pretrained(
|
||||
"runwayml/stable-diffusion-inpainting",
|
||||
custom_pipeline="img2img_inpainting",
|
||||
revision="fp16",
|
||||
torch_dtype=torch.float16,
|
||||
torch_dtype=torch.float16
|
||||
)
|
||||
pipe = pipe.to("cuda")
|
||||
|
||||
@@ -623,6 +624,8 @@ prompt = "Your prompt here!"
|
||||
image = pipe(prompt=prompt, image=init_image, inner_image=inner_image, mask_image=mask_image).images[0]
|
||||
```
|
||||
|
||||

|
||||
|
||||
### Text Based Inpainting Stable Diffusion
|
||||
|
||||
Use a text prompt to generate the mask for the area to be inpainted.
|
||||
|
||||
@@ -138,7 +138,7 @@ def ddpm_bit_scheduler_step(
|
||||
model_output: torch.FloatTensor,
|
||||
timestep: int,
|
||||
sample: torch.FloatTensor,
|
||||
predict_epsilon=True,
|
||||
prediction_type="epsilon",
|
||||
generator=None,
|
||||
return_dict: bool = True,
|
||||
) -> Union[DDPMSchedulerOutput, Tuple]:
|
||||
@@ -150,8 +150,8 @@ def ddpm_bit_scheduler_step(
|
||||
timestep (`int`): current discrete timestep in the diffusion chain.
|
||||
sample (`torch.FloatTensor`):
|
||||
current instance of sample being created by diffusion process.
|
||||
predict_epsilon (`bool`):
|
||||
optional flag to use when model predicts the samples directly instead of the noise, epsilon.
|
||||
prediction_type (`str`, default `epsilon`):
|
||||
indicates whether the model predicts the noise (epsilon), or the samples (`sample`).
|
||||
generator: random number generator.
|
||||
return_dict (`bool`): option for returning tuple rather than DDPMSchedulerOutput class
|
||||
Returns:
|
||||
@@ -174,10 +174,12 @@ def ddpm_bit_scheduler_step(
|
||||
|
||||
# 2. compute predicted original sample from predicted noise also called
|
||||
# "predicted x_0" of formula (15) from https://arxiv.org/pdf/2006.11239.pdf
|
||||
if predict_epsilon:
|
||||
if prediction_type == "epsilon":
|
||||
pred_original_sample = (sample - beta_prod_t ** (0.5) * model_output) / alpha_prod_t ** (0.5)
|
||||
else:
|
||||
elif prediction_type == "sample":
|
||||
pred_original_sample = model_output
|
||||
else:
|
||||
raise ValueError(f"Unsupported prediction_type {prediction_type}.")
|
||||
|
||||
# 3. Clip "predicted x_0"
|
||||
scale = self.bit_scale
|
||||
|
||||
@@ -78,7 +78,12 @@ class CLIPGuidedStableDiffusion(DiffusionPipeline):
|
||||
)
|
||||
|
||||
self.normalize = transforms.Normalize(mean=feature_extractor.image_mean, std=feature_extractor.image_std)
|
||||
self.make_cutouts = MakeCutouts(feature_extractor.size)
|
||||
cut_out_size = (
|
||||
feature_extractor.size
|
||||
if isinstance(feature_extractor.size, int)
|
||||
else feature_extractor.size["shortest_edge"]
|
||||
)
|
||||
self.make_cutouts = MakeCutouts(cut_out_size)
|
||||
|
||||
set_requires_grad(self.text_encoder, False)
|
||||
set_requires_grad(self.clip_model, False)
|
||||
|
||||
@@ -17,7 +17,7 @@ from diffusers.pipeline_utils import DiffusionPipeline
|
||||
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
|
||||
from diffusers.pipelines.stable_diffusion.safety_checker import StableDiffusionSafetyChecker
|
||||
from diffusers.schedulers import DDIMScheduler, LMSDiscreteScheduler, PNDMScheduler
|
||||
from diffusers.utils import logging
|
||||
from diffusers.utils import deprecate, logging
|
||||
|
||||
# TODO: remove and import from diffusers.utils when the new version of diffusers is released
|
||||
from packaging import version
|
||||
@@ -133,7 +133,7 @@ class ImagicStableDiffusionPipeline(DiffusionPipeline):
|
||||
def train(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
init_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
height: Optional[int] = 512,
|
||||
width: Optional[int] = 512,
|
||||
generator: Optional[torch.Generator] = None,
|
||||
@@ -184,6 +184,10 @@ class ImagicStableDiffusionPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
message = "Please use `image` instead of `init_image`."
|
||||
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
|
||||
image = init_image or image
|
||||
|
||||
accelerator = Accelerator(
|
||||
gradient_accumulation_steps=1,
|
||||
mixed_precision="fp16",
|
||||
@@ -241,14 +245,14 @@ class ImagicStableDiffusionPipeline(DiffusionPipeline):
|
||||
lr=embedding_learning_rate,
|
||||
)
|
||||
|
||||
if isinstance(init_image, PIL.Image.Image):
|
||||
init_image = preprocess(init_image)
|
||||
if isinstance(image, PIL.Image.Image):
|
||||
image = preprocess(image)
|
||||
|
||||
latents_dtype = text_embeddings.dtype
|
||||
init_image = init_image.to(device=self.device, dtype=latents_dtype)
|
||||
init_latent_image_dist = self.vae.encode(init_image).latent_dist
|
||||
init_image_latents = init_latent_image_dist.sample(generator=generator)
|
||||
init_image_latents = 0.18215 * init_image_latents
|
||||
image = image.to(device=self.device, dtype=latents_dtype)
|
||||
init_latent_image_dist = self.vae.encode(image).latent_dist
|
||||
image_latents = init_latent_image_dist.sample(generator=generator)
|
||||
image_latents = 0.18215 * image_latents
|
||||
|
||||
progress_bar = tqdm(range(text_embedding_optimization_steps), disable=not accelerator.is_local_main_process)
|
||||
progress_bar.set_description("Steps")
|
||||
@@ -259,12 +263,12 @@ class ImagicStableDiffusionPipeline(DiffusionPipeline):
|
||||
for _ in range(text_embedding_optimization_steps):
|
||||
with accelerator.accumulate(text_embeddings):
|
||||
# Sample noise that we'll add to the latents
|
||||
noise = torch.randn(init_image_latents.shape).to(init_image_latents.device)
|
||||
timesteps = torch.randint(1000, (1,), device=init_image_latents.device)
|
||||
noise = torch.randn(image_latents.shape).to(image_latents.device)
|
||||
timesteps = torch.randint(1000, (1,), device=image_latents.device)
|
||||
|
||||
# Add noise to the latents according to the noise magnitude at each timestep
|
||||
# (this is the forward diffusion process)
|
||||
noisy_latents = self.scheduler.add_noise(init_image_latents, noise, timesteps)
|
||||
noisy_latents = self.scheduler.add_noise(image_latents, noise, timesteps)
|
||||
|
||||
# Predict the noise residual
|
||||
noise_pred = self.unet(noisy_latents, timesteps, text_embeddings).sample
|
||||
@@ -301,12 +305,12 @@ class ImagicStableDiffusionPipeline(DiffusionPipeline):
|
||||
for _ in range(model_fine_tuning_optimization_steps):
|
||||
with accelerator.accumulate(self.unet.parameters()):
|
||||
# Sample noise that we'll add to the latents
|
||||
noise = torch.randn(init_image_latents.shape).to(init_image_latents.device)
|
||||
timesteps = torch.randint(1000, (1,), device=init_image_latents.device)
|
||||
noise = torch.randn(image_latents.shape).to(image_latents.device)
|
||||
timesteps = torch.randint(1000, (1,), device=image_latents.device)
|
||||
|
||||
# Add noise to the latents according to the noise magnitude at each timestep
|
||||
# (this is the forward diffusion process)
|
||||
noisy_latents = self.scheduler.add_noise(init_image_latents, noise, timesteps)
|
||||
noisy_latents = self.scheduler.add_noise(image_latents, noise, timesteps)
|
||||
|
||||
# Predict the noise residual
|
||||
noise_pred = self.unet(noisy_latents, timesteps, text_embeddings).sample
|
||||
|
||||
@@ -555,7 +555,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
init_image: Union[torch.FloatTensor, PIL.Image.Image] = None,
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image] = None,
|
||||
mask_image: Union[torch.FloatTensor, PIL.Image.Image] = None,
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
@@ -583,11 +583,11 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
negative_prompt (`str` or `List[str]`, *optional*):
|
||||
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
|
||||
if `guidance_scale` is less than `1`).
|
||||
init_image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process.
|
||||
mask_image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, to mask `init_image`. White pixels in the mask will be
|
||||
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be
|
||||
replaced by noise and therefore repainted, while black pixels will be preserved. If `mask_image` is a
|
||||
PIL image, it will be converted to a single channel (luminance) before use. If it's a tensor, it should
|
||||
contain one color channel (L) instead of 3, so the expected shape would be `(B, H, W, 1)`.
|
||||
@@ -605,11 +605,11 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`,
|
||||
usually at the expense of lower image quality.
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
|
||||
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1.
|
||||
`image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
|
||||
noise will be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
|
||||
num_images_per_prompt (`int`, *optional*, defaults to 1):
|
||||
The number of images to generate per prompt.
|
||||
eta (`float`, *optional*, defaults to 0.0):
|
||||
@@ -648,6 +648,9 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
message = "Please use `image` instead of `init_image`."
|
||||
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
|
||||
image = init_image or image
|
||||
|
||||
if isinstance(prompt, str):
|
||||
batch_size = 1
|
||||
@@ -714,7 +717,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
mask = None
|
||||
noise = None
|
||||
|
||||
if init_image is None:
|
||||
if image is None:
|
||||
# get the initial random noise unless the user supplied it
|
||||
|
||||
# Unlike in other pipelines, latents need to be generated in the target device
|
||||
@@ -753,11 +756,11 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
# scale the initial noise by the standard deviation required by the scheduler
|
||||
latents = latents * self.scheduler.init_noise_sigma
|
||||
else:
|
||||
if isinstance(init_image, PIL.Image.Image):
|
||||
init_image = preprocess_image(init_image)
|
||||
if isinstance(image, PIL.Image.Image):
|
||||
image = preprocess_image(image)
|
||||
# encode the init image into latents and scale the latents
|
||||
init_image = init_image.to(device=self.device, dtype=latents_dtype)
|
||||
init_latent_dist = self.vae.encode(init_image).latent_dist
|
||||
image = image.to(device=self.device, dtype=latents_dtype)
|
||||
init_latent_dist = self.vae.encode(image).latent_dist
|
||||
init_latents = init_latent_dist.sample(generator=generator)
|
||||
init_latents = 0.18215 * init_latents
|
||||
init_latents = torch.cat([init_latents] * batch_size * num_images_per_prompt, dim=0)
|
||||
@@ -772,7 +775,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
|
||||
# check sizes
|
||||
if not mask.shape == init_latents.shape:
|
||||
raise ValueError("The mask and init_image should be the same size!")
|
||||
raise ValueError("The mask and image should be the same size!")
|
||||
|
||||
# get the original timestep using init_timestep
|
||||
offset = self.scheduler.config.get("steps_offset", 0)
|
||||
@@ -961,7 +964,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
|
||||
def img2img(
|
||||
self,
|
||||
init_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
prompt: Union[str, List[str]],
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
strength: float = 0.8,
|
||||
@@ -980,7 +983,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
r"""
|
||||
Function for image-to-image generation.
|
||||
Args:
|
||||
init_image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process.
|
||||
prompt (`str` or `List[str]`):
|
||||
@@ -989,11 +992,11 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
|
||||
if `guidance_scale` is less than `1`).
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
|
||||
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1.
|
||||
`image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
|
||||
noise will be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
expense of slower inference. This parameter will be modulated by `strength`.
|
||||
@@ -1035,7 +1038,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
return self.__call__(
|
||||
prompt=prompt,
|
||||
negative_prompt=negative_prompt,
|
||||
init_image=init_image,
|
||||
image=image,
|
||||
num_inference_steps=num_inference_steps,
|
||||
guidance_scale=guidance_scale,
|
||||
strength=strength,
|
||||
@@ -1052,7 +1055,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
|
||||
def inpaint(
|
||||
self,
|
||||
init_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
mask_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
prompt: Union[str, List[str]],
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -1072,11 +1075,11 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
r"""
|
||||
Function for inpaint.
|
||||
Args:
|
||||
init_image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process. This is the image whose masked region will be inpainted.
|
||||
mask_image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, to mask `init_image`. White pixels in the mask will be
|
||||
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be
|
||||
replaced by noise and therefore repainted, while black pixels will be preserved. If `mask_image` is a
|
||||
PIL image, it will be converted to a single channel (luminance) before use. If it's a tensor, it should
|
||||
contain one color channel (L) instead of 3, so the expected shape would be `(B, H, W, 1)`.
|
||||
@@ -1088,7 +1091,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to inpaint the masked area. Must be between 0 and 1. When `strength`
|
||||
is 1, the denoising process will be run on the masked area for the full number of iterations specified
|
||||
in `num_inference_steps`. `init_image` will be used as a reference for the masked area, adding more
|
||||
in `num_inference_steps`. `image` will be used as a reference for the masked area, adding more
|
||||
noise to that region the larger the `strength`. If `strength` is 0, no inpainting will occur.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The reference number of denoising steps. More denoising steps usually lead to a higher quality image at
|
||||
@@ -1131,7 +1134,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
return self.__call__(
|
||||
prompt=prompt,
|
||||
negative_prompt=negative_prompt,
|
||||
init_image=init_image,
|
||||
image=image,
|
||||
mask_image=mask_image,
|
||||
num_inference_steps=num_inference_steps,
|
||||
guidance_scale=guidance_scale,
|
||||
|
||||
@@ -10,7 +10,7 @@ from diffusers.onnx_utils import OnnxRuntimeModel
|
||||
from diffusers.pipeline_utils import DiffusionPipeline
|
||||
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
|
||||
from diffusers.schedulers import DDIMScheduler, LMSDiscreteScheduler, PNDMScheduler
|
||||
from diffusers.utils import logging
|
||||
from diffusers.utils import deprecate, logging
|
||||
|
||||
# TODO: remove and import from diffusers.utils when the new version of diffusers is released
|
||||
from packaging import version
|
||||
@@ -441,7 +441,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
init_image: Union[np.ndarray, PIL.Image.Image] = None,
|
||||
image: Union[np.ndarray, PIL.Image.Image] = None,
|
||||
mask_image: Union[np.ndarray, PIL.Image.Image] = None,
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
@@ -469,11 +469,11 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
negative_prompt (`str` or `List[str]`, *optional*):
|
||||
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
|
||||
if `guidance_scale` is less than `1`).
|
||||
init_image (`np.ndarray` or `PIL.Image.Image`):
|
||||
image (`np.ndarray` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process.
|
||||
mask_image (`np.ndarray` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, to mask `init_image`. White pixels in the mask will be
|
||||
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be
|
||||
replaced by noise and therefore repainted, while black pixels will be preserved. If `mask_image` is a
|
||||
PIL image, it will be converted to a single channel (luminance) before use. If it's a tensor, it should
|
||||
contain one color channel (L) instead of 3, so the expected shape would be `(B, H, W, 1)`.
|
||||
@@ -491,11 +491,11 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`,
|
||||
usually at the expense of lower image quality.
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
|
||||
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1.
|
||||
`image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
|
||||
noise will be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
|
||||
num_images_per_prompt (`int`, *optional*, defaults to 1):
|
||||
The number of images to generate per prompt.
|
||||
eta (`float`, *optional*, defaults to 0.0):
|
||||
@@ -533,6 +533,9 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
message = "Please use `image` instead of `init_image`."
|
||||
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
|
||||
image = init_image or image
|
||||
|
||||
if isinstance(prompt, str):
|
||||
batch_size = 1
|
||||
@@ -598,7 +601,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
mask = None
|
||||
noise = None
|
||||
|
||||
if init_image is None:
|
||||
if image is None:
|
||||
latents_shape = (
|
||||
batch_size * num_images_per_prompt,
|
||||
4,
|
||||
@@ -616,11 +619,11 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
# scale the initial noise by the standard deviation required by the scheduler
|
||||
latents = latents * self.scheduler.init_noise_sigma
|
||||
else:
|
||||
if isinstance(init_image, PIL.Image.Image):
|
||||
init_image = preprocess_image(init_image)
|
||||
if isinstance(image, PIL.Image.Image):
|
||||
image = preprocess_image(image)
|
||||
# encode the init image into latents and scale the latents
|
||||
init_image = init_image.astype(latents_dtype)
|
||||
init_latents = self.vae_encoder(sample=init_image)[0]
|
||||
image = image.astype(latents_dtype)
|
||||
init_latents = self.vae_encoder(sample=image)[0]
|
||||
init_latents = 0.18215 * init_latents
|
||||
init_latents = np.concatenate([init_latents] * batch_size * num_images_per_prompt)
|
||||
init_latents_orig = init_latents
|
||||
@@ -635,7 +638,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
# check sizes
|
||||
if not mask.shape == init_latents.shape:
|
||||
print(mask.shape, init_latents.shape)
|
||||
raise ValueError("The mask and init_image should be the same size!")
|
||||
raise ValueError("The mask and image should be the same size!")
|
||||
|
||||
# get the original timestep using init_timestep
|
||||
offset = self.scheduler.config.get("steps_offset", 0)
|
||||
@@ -828,7 +831,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
|
||||
def img2img(
|
||||
self,
|
||||
init_image: Union[np.ndarray, PIL.Image.Image],
|
||||
image: Union[np.ndarray, PIL.Image.Image],
|
||||
prompt: Union[str, List[str]],
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
strength: float = 0.8,
|
||||
@@ -847,7 +850,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
r"""
|
||||
Function for image-to-image generation.
|
||||
Args:
|
||||
init_image (`np.ndarray` or `PIL.Image.Image`):
|
||||
image (`np.ndarray` or `PIL.Image.Image`):
|
||||
`Image`, or ndarray representing an image batch, that will be used as the starting point for the
|
||||
process.
|
||||
prompt (`str` or `List[str]`):
|
||||
@@ -856,11 +859,11 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
|
||||
if `guidance_scale` is less than `1`).
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
|
||||
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1.
|
||||
`image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
|
||||
noise will be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
expense of slower inference. This parameter will be modulated by `strength`.
|
||||
@@ -901,7 +904,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
return self.__call__(
|
||||
prompt=prompt,
|
||||
negative_prompt=negative_prompt,
|
||||
init_image=init_image,
|
||||
image=image,
|
||||
num_inference_steps=num_inference_steps,
|
||||
guidance_scale=guidance_scale,
|
||||
strength=strength,
|
||||
@@ -918,7 +921,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
|
||||
def inpaint(
|
||||
self,
|
||||
init_image: Union[np.ndarray, PIL.Image.Image],
|
||||
image: Union[np.ndarray, PIL.Image.Image],
|
||||
mask_image: Union[np.ndarray, PIL.Image.Image],
|
||||
prompt: Union[str, List[str]],
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -938,11 +941,11 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
r"""
|
||||
Function for inpaint.
|
||||
Args:
|
||||
init_image (`np.ndarray` or `PIL.Image.Image`):
|
||||
image (`np.ndarray` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process. This is the image whose masked region will be inpainted.
|
||||
mask_image (`np.ndarray` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, to mask `init_image`. White pixels in the mask will be
|
||||
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be
|
||||
replaced by noise and therefore repainted, while black pixels will be preserved. If `mask_image` is a
|
||||
PIL image, it will be converted to a single channel (luminance) before use. If it's a tensor, it should
|
||||
contain one color channel (L) instead of 3, so the expected shape would be `(B, H, W, 1)`.
|
||||
@@ -954,7 +957,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to inpaint the masked area. Must be between 0 and 1. When `strength`
|
||||
is 1, the denoising process will be run on the masked area for the full number of iterations specified
|
||||
in `num_inference_steps`. `init_image` will be used as a reference for the masked area, adding more
|
||||
in `num_inference_steps`. `image` will be used as a reference for the masked area, adding more
|
||||
noise to that region the larger the `strength`. If `strength` is 0, no inpainting will occur.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The reference number of denoising steps. More denoising steps usually lead to a higher quality image at
|
||||
@@ -996,7 +999,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
|
||||
return self.__call__(
|
||||
prompt=prompt,
|
||||
negative_prompt=negative_prompt,
|
||||
init_image=init_image,
|
||||
image=image,
|
||||
mask_image=mask_image,
|
||||
num_inference_steps=num_inference_steps,
|
||||
guidance_scale=guidance_scale,
|
||||
|
||||
@@ -121,7 +121,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
|
||||
def inpaint(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
init_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
mask_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
strength: float = 0.8,
|
||||
num_inference_steps: Optional[int] = 50,
|
||||
@@ -138,7 +138,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
|
||||
# For more information on how this function works, please see: https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion#diffusers.StableDiffusionImg2ImgPipeline
|
||||
return StableDiffusionInpaintPipelineLegacy(**self.components)(
|
||||
prompt=prompt,
|
||||
init_image=init_image,
|
||||
image=image,
|
||||
mask_image=mask_image,
|
||||
strength=strength,
|
||||
num_inference_steps=num_inference_steps,
|
||||
@@ -156,7 +156,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
|
||||
def img2img(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
init_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
strength: float = 0.8,
|
||||
num_inference_steps: Optional[int] = 50,
|
||||
guidance_scale: Optional[float] = 7.5,
|
||||
@@ -173,7 +173,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
|
||||
# For more information on how this function works, please see: https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion#diffusers.StableDiffusionImg2ImgPipeline
|
||||
return StableDiffusionImg2ImgPipeline(**self.components)(
|
||||
prompt=prompt,
|
||||
init_image=init_image,
|
||||
image=image,
|
||||
strength=strength,
|
||||
num_inference_steps=num_inference_steps,
|
||||
guidance_scale=guidance_scale,
|
||||
|
||||
@@ -39,6 +39,8 @@ Now let's get our dataset. Download images from [here](https://drive.google.com/
|
||||
|
||||
And launch the training using
|
||||
|
||||
**___Note: Change the `resolution` to 768 if you are using the [stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2) 768x768 model.___**
|
||||
|
||||
```bash
|
||||
export MODEL_NAME="CompVis/stable-diffusion-v1-4"
|
||||
export INSTANCE_DIR="path-to-instance-images"
|
||||
@@ -193,6 +195,17 @@ accelerate launch train_dreambooth.py \
|
||||
--max_train_steps=800
|
||||
```
|
||||
|
||||
### Using DreamBooth for other pipelines than Stable Diffusion
|
||||
|
||||
Altdiffusion also support dreambooth now, the runing comman is basically the same as abouve, all you need to do is replace the `MODEL_NAME` like this:
|
||||
One can now simply change the `pretrained_model_name_or_path` to another architecture such as [`AltDiffusion`](https://huggingface.co/docs/diffusers/api/pipelines/alt_diffusion).
|
||||
|
||||
```
|
||||
export MODEL_NAME="CompVis/stable-diffusion-v1-4" --> export MODEL_NAME="BAAI/AltDiffusion-m9"
|
||||
or
|
||||
export MODEL_NAME="CompVis/stable-diffusion-v1-4" --> export MODEL_NAME="BAAI/AltDiffusion"
|
||||
```
|
||||
|
||||
### Inference
|
||||
|
||||
Once you have trained a model using above command, the inference can be done simply using the `StableDiffusionPipeline`. Make sure to include the `identifier`(e.g. sks in above example) in your prompt.
|
||||
|
||||
@@ -14,18 +14,38 @@ from torch.utils.data import Dataset
|
||||
from accelerate import Accelerator
|
||||
from accelerate.logging import get_logger
|
||||
from accelerate.utils import set_seed
|
||||
from diffusers import AutoencoderKL, DDPMScheduler, StableDiffusionPipeline, UNet2DConditionModel
|
||||
from diffusers import AutoencoderKL, DDPMScheduler, DiffusionPipeline, UNet2DConditionModel
|
||||
from diffusers.optimization import get_scheduler
|
||||
from huggingface_hub import HfFolder, Repository, whoami
|
||||
from PIL import Image
|
||||
from torchvision import transforms
|
||||
from tqdm.auto import tqdm
|
||||
from transformers import CLIPTextModel, CLIPTokenizer
|
||||
from transformers import AutoTokenizer, PretrainedConfig
|
||||
|
||||
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
def import_model_class_from_model_name_or_path(pretrained_model_name_or_path: str):
|
||||
text_encoder_config = PretrainedConfig.from_pretrained(
|
||||
pretrained_model_name_or_path,
|
||||
subfolder="text_encoder",
|
||||
revision=args.revision,
|
||||
)
|
||||
model_class = text_encoder_config.architectures[0]
|
||||
|
||||
if model_class == "CLIPTextModel":
|
||||
from transformers import CLIPTextModel
|
||||
|
||||
return CLIPTextModel
|
||||
elif model_class == "RobertaSeriesModelWithTransformation":
|
||||
from diffusers.pipelines.alt_diffusion.modeling_roberta_series import RobertaSeriesModelWithTransformation
|
||||
|
||||
return RobertaSeriesModelWithTransformation
|
||||
else:
|
||||
raise ValueError(f"{model_class} is not supported.")
|
||||
|
||||
|
||||
def parse_args(input_args=None):
|
||||
parser = argparse.ArgumentParser(description="Simple example of a training script.")
|
||||
parser.add_argument(
|
||||
@@ -124,6 +144,7 @@ def parse_args(input_args=None):
|
||||
default=None,
|
||||
help="Total number of training steps to perform. If provided, overrides num_train_epochs.",
|
||||
)
|
||||
parser.add_argument("--save_steps", type=int, default=500, help="Save checkpoint every X updates steps.")
|
||||
parser.add_argument(
|
||||
"--gradient_accumulation_steps",
|
||||
type=int,
|
||||
@@ -356,7 +377,7 @@ def main(args):
|
||||
|
||||
if cur_class_images < args.num_class_images:
|
||||
torch_dtype = torch.float16 if accelerator.device.type == "cuda" else torch.float32
|
||||
pipeline = StableDiffusionPipeline.from_pretrained(
|
||||
pipeline = DiffusionPipeline.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
torch_dtype=torch_dtype,
|
||||
safety_checker=None,
|
||||
@@ -406,19 +427,24 @@ def main(args):
|
||||
|
||||
# Load the tokenizer
|
||||
if args.tokenizer_name:
|
||||
tokenizer = CLIPTokenizer.from_pretrained(
|
||||
tokenizer = AutoTokenizer.from_pretrained(
|
||||
args.tokenizer_name,
|
||||
revision=args.revision,
|
||||
use_fast=False,
|
||||
)
|
||||
elif args.pretrained_model_name_or_path:
|
||||
tokenizer = CLIPTokenizer.from_pretrained(
|
||||
tokenizer = AutoTokenizer.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
subfolder="tokenizer",
|
||||
revision=args.revision,
|
||||
use_fast=False,
|
||||
)
|
||||
|
||||
# import correct text encoder class
|
||||
text_encoder_cls = import_model_class_from_model_name_or_path(args.pretrained_model_name_or_path)
|
||||
|
||||
# Load models and create wrapper for stable diffusion
|
||||
text_encoder = CLIPTextModel.from_pretrained(
|
||||
text_encoder = text_encoder_cls.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
subfolder="text_encoder",
|
||||
revision=args.revision,
|
||||
@@ -603,23 +629,31 @@ def main(args):
|
||||
encoder_hidden_states = text_encoder(batch["input_ids"])[0]
|
||||
|
||||
# Predict the noise residual
|
||||
noise_pred = unet(noisy_latents, timesteps, encoder_hidden_states).sample
|
||||
model_pred = unet(noisy_latents, timesteps, encoder_hidden_states).sample
|
||||
|
||||
# Get the target for loss depending on the prediction type
|
||||
if noise_scheduler.config.prediction_type == "epsilon":
|
||||
target = noise
|
||||
elif noise_scheduler.config.prediction_type == "v_prediction":
|
||||
target = noise_scheduler.get_velocity(latents, noise, timesteps)
|
||||
else:
|
||||
raise ValueError(f"Unknown prediction type {noise_scheduler.config.prediction_type}")
|
||||
|
||||
if args.with_prior_preservation:
|
||||
# Chunk the noise and noise_pred into two parts and compute the loss on each part separately.
|
||||
noise_pred, noise_pred_prior = torch.chunk(noise_pred, 2, dim=0)
|
||||
noise, noise_prior = torch.chunk(noise, 2, dim=0)
|
||||
# Chunk the noise and model_pred into two parts and compute the loss on each part separately.
|
||||
model_pred, model_pred_prior = torch.chunk(model_pred, 2, dim=0)
|
||||
target, target_prior = torch.chunk(target, 2, dim=0)
|
||||
|
||||
# Compute instance loss
|
||||
loss = F.mse_loss(noise_pred.float(), noise.float(), reduction="none").mean([1, 2, 3]).mean()
|
||||
loss = F.mse_loss(model_pred.float(), target.float(), reduction="none").mean([1, 2, 3]).mean()
|
||||
|
||||
# Compute prior loss
|
||||
prior_loss = F.mse_loss(noise_pred_prior.float(), noise_prior.float(), reduction="mean")
|
||||
prior_loss = F.mse_loss(model_pred_prior.float(), target_prior.float(), reduction="mean")
|
||||
|
||||
# Add the prior loss to the instance loss.
|
||||
loss = loss + args.prior_loss_weight * prior_loss
|
||||
else:
|
||||
loss = F.mse_loss(noise_pred.float(), noise.float(), reduction="mean")
|
||||
loss = F.mse_loss(model_pred.float(), target.float(), reduction="mean")
|
||||
|
||||
accelerator.backward(loss)
|
||||
if accelerator.sync_gradients:
|
||||
@@ -638,6 +672,17 @@ def main(args):
|
||||
progress_bar.update(1)
|
||||
global_step += 1
|
||||
|
||||
if global_step % args.save_steps == 0:
|
||||
if accelerator.is_main_process:
|
||||
pipeline = DiffusionPipeline.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
unet=accelerator.unwrap_model(unet),
|
||||
text_encoder=accelerator.unwrap_model(text_encoder),
|
||||
revision=args.revision,
|
||||
)
|
||||
save_path = os.path.join(args.output_dir, f"checkpoint-{global_step}")
|
||||
pipeline.save_pretrained(save_path)
|
||||
|
||||
logs = {"loss": loss.detach().item(), "lr": lr_scheduler.get_last_lr()[0]}
|
||||
progress_bar.set_postfix(**logs)
|
||||
accelerator.log(logs, step=global_step)
|
||||
@@ -649,7 +694,7 @@ def main(args):
|
||||
|
||||
# Create the pipeline using using the trained modules and save it.
|
||||
if accelerator.is_main_process:
|
||||
pipeline = StableDiffusionPipeline.from_pretrained(
|
||||
pipeline = DiffusionPipeline.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
unet=accelerator.unwrap_model(unet),
|
||||
text_encoder=accelerator.unwrap_model(text_encoder),
|
||||
|
||||
@@ -42,6 +42,8 @@ If you have already cloned the repo, then you won't need to go through these ste
|
||||
#### Hardware
|
||||
With `gradient_checkpointing` and `mixed_precision` it should be possible to fine tune the model on a single 24GB GPU. For higher `batch_size` and faster training it's better to use GPUs with >30GB memory.
|
||||
|
||||
**___Note: Change the `resolution` to 768 if you are using the [stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2) 768x768 model.___**
|
||||
|
||||
```bash
|
||||
export MODEL_NAME="CompVis/stable-diffusion-v1-4"
|
||||
export dataset_name="lambdalabs/pokemon-blip-captions"
|
||||
|
||||
@@ -15,13 +15,12 @@ from accelerate import Accelerator
|
||||
from accelerate.logging import get_logger
|
||||
from accelerate.utils import set_seed
|
||||
from datasets import load_dataset
|
||||
from diffusers import AutoencoderKL, DDPMScheduler, PNDMScheduler, StableDiffusionPipeline, UNet2DConditionModel
|
||||
from diffusers import AutoencoderKL, DDPMScheduler, StableDiffusionPipeline, UNet2DConditionModel
|
||||
from diffusers.optimization import get_scheduler
|
||||
from diffusers.pipelines.stable_diffusion import StableDiffusionSafetyChecker
|
||||
from huggingface_hub import HfFolder, Repository, whoami
|
||||
from torchvision import transforms
|
||||
from tqdm.auto import tqdm
|
||||
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
|
||||
from transformers import CLIPTextModel, CLIPTokenizer
|
||||
|
||||
|
||||
logger = get_logger(__name__)
|
||||
@@ -36,6 +35,13 @@ def parse_args():
|
||||
required=True,
|
||||
help="Path to pretrained model or model identifier from huggingface.co/models.",
|
||||
)
|
||||
parser.add_argument(
|
||||
"--revision",
|
||||
type=str,
|
||||
default=None,
|
||||
required=False,
|
||||
help="Revision of pretrained model identifier from huggingface.co/models.",
|
||||
)
|
||||
parser.add_argument(
|
||||
"--dataset_name",
|
||||
type=str,
|
||||
@@ -335,10 +341,24 @@ def main():
|
||||
os.makedirs(args.output_dir, exist_ok=True)
|
||||
|
||||
# Load models and create wrapper for stable diffusion
|
||||
tokenizer = CLIPTokenizer.from_pretrained(args.pretrained_model_name_or_path, subfolder="tokenizer")
|
||||
text_encoder = CLIPTextModel.from_pretrained(args.pretrained_model_name_or_path, subfolder="text_encoder")
|
||||
vae = AutoencoderKL.from_pretrained(args.pretrained_model_name_or_path, subfolder="vae")
|
||||
unet = UNet2DConditionModel.from_pretrained(args.pretrained_model_name_or_path, subfolder="unet")
|
||||
tokenizer = CLIPTokenizer.from_pretrained(
|
||||
args.pretrained_model_name_or_path, subfolder="tokenizer", revision=args.revision
|
||||
)
|
||||
text_encoder = CLIPTextModel.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
subfolder="text_encoder",
|
||||
revision=args.revision,
|
||||
)
|
||||
vae = AutoencoderKL.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
subfolder="vae",
|
||||
revision=args.revision,
|
||||
)
|
||||
unet = UNet2DConditionModel.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
subfolder="unet",
|
||||
revision=args.revision,
|
||||
)
|
||||
|
||||
# Freeze vae and text_encoder
|
||||
vae.requires_grad_(False)
|
||||
@@ -562,9 +582,17 @@ def main():
|
||||
# Get the text embedding for conditioning
|
||||
encoder_hidden_states = text_encoder(batch["input_ids"])[0]
|
||||
|
||||
# Get the target for loss depending on the prediction type
|
||||
if noise_scheduler.config.prediction_type == "epsilon":
|
||||
target = noise
|
||||
elif noise_scheduler.config.prediction_type == "v_prediction":
|
||||
target = noise_scheduler.get_velocity(latents, noise, timesteps)
|
||||
else:
|
||||
raise ValueError(f"Unknown prediction type {noise_scheduler.config.prediction_type}")
|
||||
|
||||
# Predict the noise residual and compute loss
|
||||
noise_pred = unet(noisy_latents, timesteps, encoder_hidden_states).sample
|
||||
loss = F.mse_loss(noise_pred.float(), noise.float(), reduction="mean")
|
||||
model_pred = unet(noisy_latents, timesteps, encoder_hidden_states).sample
|
||||
loss = F.mse_loss(model_pred.float(), target.float(), reduction="mean")
|
||||
|
||||
# Gather the losses across all processes for logging (if we use distributed training).
|
||||
avg_loss = accelerator.gather(loss.repeat(args.train_batch_size)).mean()
|
||||
@@ -600,14 +628,12 @@ def main():
|
||||
if args.use_ema:
|
||||
ema_unet.copy_to(unet.parameters())
|
||||
|
||||
pipeline = StableDiffusionPipeline(
|
||||
pipeline = StableDiffusionPipeline.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
text_encoder=text_encoder,
|
||||
vae=vae,
|
||||
unet=unet,
|
||||
tokenizer=tokenizer,
|
||||
scheduler=PNDMScheduler.from_config(args.pretrained_model_name_or_path, subfolder="scheduler"),
|
||||
safety_checker=StableDiffusionSafetyChecker.from_pretrained("CompVis/stable-diffusion-safety-checker"),
|
||||
feature_extractor=CLIPFeatureExtractor.from_pretrained("openai/clip-vit-base-patch32"),
|
||||
revision=args.revision,
|
||||
)
|
||||
pipeline.save_pretrained(args.output_dir)
|
||||
|
||||
|
||||
@@ -47,6 +47,8 @@ Now let's get our dataset.Download 3-4 images from [here](https://drive.google.c
|
||||
|
||||
And launch the training using
|
||||
|
||||
**___Note: Change the `resolution` to 768 if you are using the [stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2) 768x768 model.___**
|
||||
|
||||
```bash
|
||||
export MODEL_NAME="runwayml/stable-diffusion-v1-5"
|
||||
export DATA_DIR="path-to-dir-containing-images"
|
||||
|
||||
@@ -16,9 +16,8 @@ import PIL
|
||||
from accelerate import Accelerator
|
||||
from accelerate.logging import get_logger
|
||||
from accelerate.utils import set_seed
|
||||
from diffusers import AutoencoderKL, DDPMScheduler, PNDMScheduler, StableDiffusionPipeline, UNet2DConditionModel
|
||||
from diffusers import AutoencoderKL, DDPMScheduler, StableDiffusionPipeline, UNet2DConditionModel
|
||||
from diffusers.optimization import get_scheduler
|
||||
from diffusers.pipelines.stable_diffusion import StableDiffusionSafetyChecker
|
||||
from huggingface_hub import HfFolder, Repository, whoami
|
||||
|
||||
# TODO: remove and import from diffusers.utils when the new version of diffusers is released
|
||||
@@ -26,7 +25,7 @@ from packaging import version
|
||||
from PIL import Image
|
||||
from torchvision import transforms
|
||||
from tqdm.auto import tqdm
|
||||
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
|
||||
from transformers import CLIPTextModel, CLIPTokenizer
|
||||
|
||||
|
||||
if version.parse(version.parse(PIL.__version__).base_version) >= version.parse("9.1.0"):
|
||||
@@ -51,11 +50,11 @@ else:
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
def save_progress(text_encoder, placeholder_token_id, accelerator, args):
|
||||
def save_progress(text_encoder, placeholder_token_id, accelerator, args, save_path):
|
||||
logger.info("Saving embeddings")
|
||||
learned_embeds = accelerator.unwrap_model(text_encoder).get_input_embeddings().weight[placeholder_token_id]
|
||||
learned_embeds_dict = {args.placeholder_token: learned_embeds.detach().cpu()}
|
||||
torch.save(learned_embeds_dict, os.path.join(args.output_dir, "learned_embeds.bin"))
|
||||
torch.save(learned_embeds_dict, save_path)
|
||||
|
||||
|
||||
def parse_args():
|
||||
@@ -73,6 +72,13 @@ def parse_args():
|
||||
required=True,
|
||||
help="Path to pretrained model or model identifier from huggingface.co/models.",
|
||||
)
|
||||
parser.add_argument(
|
||||
"--revision",
|
||||
type=str,
|
||||
default=None,
|
||||
required=False,
|
||||
help="Revision of pretrained model identifier from huggingface.co/models.",
|
||||
)
|
||||
parser.add_argument(
|
||||
"--tokenizer_name",
|
||||
type=str,
|
||||
@@ -405,9 +411,21 @@ def main():
|
||||
placeholder_token_id = tokenizer.convert_tokens_to_ids(args.placeholder_token)
|
||||
|
||||
# Load models and create wrapper for stable diffusion
|
||||
text_encoder = CLIPTextModel.from_pretrained(args.pretrained_model_name_or_path, subfolder="text_encoder")
|
||||
vae = AutoencoderKL.from_pretrained(args.pretrained_model_name_or_path, subfolder="vae")
|
||||
unet = UNet2DConditionModel.from_pretrained(args.pretrained_model_name_or_path, subfolder="unet")
|
||||
text_encoder = CLIPTextModel.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
subfolder="text_encoder",
|
||||
revision=args.revision,
|
||||
)
|
||||
vae = AutoencoderKL.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
subfolder="vae",
|
||||
revision=args.revision,
|
||||
)
|
||||
unet = UNet2DConditionModel.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
subfolder="unet",
|
||||
revision=args.revision,
|
||||
)
|
||||
|
||||
# Resize the token embeddings as we are adding new special tokens to the tokenizer
|
||||
text_encoder.resize_token_embeddings(len(tokenizer))
|
||||
@@ -532,9 +550,17 @@ def main():
|
||||
encoder_hidden_states = text_encoder(batch["input_ids"])[0]
|
||||
|
||||
# Predict the noise residual
|
||||
noise_pred = unet(noisy_latents, timesteps, encoder_hidden_states).sample
|
||||
model_pred = unet(noisy_latents, timesteps, encoder_hidden_states).sample
|
||||
|
||||
loss = F.mse_loss(noise_pred, noise, reduction="none").mean([1, 2, 3]).mean()
|
||||
# Get the target for loss depending on the prediction type
|
||||
if noise_scheduler.config.prediction_type == "epsilon":
|
||||
target = noise
|
||||
elif noise_scheduler.config.prediction_type == "v_prediction":
|
||||
target = noise_scheduler.get_velocity(latents, noise, timesteps)
|
||||
else:
|
||||
raise ValueError(f"Unknown prediction type {noise_scheduler.config.prediction_type}")
|
||||
|
||||
loss = F.mse_loss(model_pred, target, reduction="none").mean([1, 2, 3]).mean()
|
||||
accelerator.backward(loss)
|
||||
|
||||
# Zero out the gradients for all token embeddings except the newly added
|
||||
@@ -556,7 +582,8 @@ def main():
|
||||
progress_bar.update(1)
|
||||
global_step += 1
|
||||
if global_step % args.save_steps == 0:
|
||||
save_progress(text_encoder, placeholder_token_id, accelerator, args)
|
||||
save_path = os.path.join(args.output_dir, f"learned_embeds-steps-{global_step}.bin")
|
||||
save_progress(text_encoder, placeholder_token_id, accelerator, args, save_path)
|
||||
|
||||
logs = {"loss": loss.detach().item(), "lr": lr_scheduler.get_last_lr()[0]}
|
||||
progress_bar.set_postfix(**logs)
|
||||
@@ -569,18 +596,18 @@ def main():
|
||||
|
||||
# Create the pipeline using using the trained modules and save it.
|
||||
if accelerator.is_main_process:
|
||||
pipeline = StableDiffusionPipeline(
|
||||
pipeline = StableDiffusionPipeline.from_pretrained(
|
||||
args.pretrained_model_name_or_path,
|
||||
text_encoder=accelerator.unwrap_model(text_encoder),
|
||||
tokenizer=tokenizer,
|
||||
vae=vae,
|
||||
unet=unet,
|
||||
tokenizer=tokenizer,
|
||||
scheduler=PNDMScheduler.from_config(args.pretrained_model_name_or_path, subfolder="scheduler"),
|
||||
safety_checker=StableDiffusionSafetyChecker.from_pretrained("CompVis/stable-diffusion-safety-checker"),
|
||||
feature_extractor=CLIPFeatureExtractor.from_pretrained("openai/clip-vit-base-patch32"),
|
||||
revision=args.revision,
|
||||
)
|
||||
pipeline.save_pretrained(args.output_dir)
|
||||
# Also save the newly trained embeddings
|
||||
save_progress(text_encoder, placeholder_token_id, accelerator, args)
|
||||
save_path = os.path.join(args.output_dir, "learned_embeds.bin")
|
||||
save_progress(text_encoder, placeholder_token_id, accelerator, args, save_path)
|
||||
|
||||
if args.push_to_hub:
|
||||
repo.push_to_hub(commit_message="End of training", blocking=False, auto_lfs_prune=True)
|
||||
|
||||
@@ -14,7 +14,6 @@ from datasets import load_dataset
|
||||
from diffusers import DDPMPipeline, DDPMScheduler, UNet2DModel, __version__
|
||||
from diffusers.optimization import get_scheduler
|
||||
from diffusers.training_utils import EMAModel
|
||||
from diffusers.utils import deprecate
|
||||
from huggingface_hub import HfFolder, Repository, whoami
|
||||
from packaging import version
|
||||
from torchvision.transforms import (
|
||||
@@ -194,9 +193,10 @@ def parse_args():
|
||||
)
|
||||
|
||||
parser.add_argument(
|
||||
"--predict_epsilon",
|
||||
action="store_true",
|
||||
default=True,
|
||||
"--prediction_type",
|
||||
type=str,
|
||||
default="epsilon",
|
||||
choices=["epsilon", "sample"],
|
||||
help="Whether the model should predict the 'epsilon'/noise error or directly the reconstructed image 'x0'.",
|
||||
)
|
||||
|
||||
@@ -256,13 +256,13 @@ def main(args):
|
||||
"UpBlock2D",
|
||||
),
|
||||
)
|
||||
accepts_predict_epsilon = "predict_epsilon" in set(inspect.signature(DDPMScheduler.__init__).parameters.keys())
|
||||
accepts_prediction_type = "prediction_type" in set(inspect.signature(DDPMScheduler.__init__).parameters.keys())
|
||||
|
||||
if accepts_predict_epsilon:
|
||||
if accepts_prediction_type:
|
||||
noise_scheduler = DDPMScheduler(
|
||||
num_train_timesteps=args.ddpm_num_steps,
|
||||
beta_schedule=args.ddpm_beta_schedule,
|
||||
predict_epsilon=args.predict_epsilon,
|
||||
prediction_type=args.prediction_type,
|
||||
)
|
||||
else:
|
||||
noise_scheduler = DDPMScheduler(num_train_timesteps=args.ddpm_num_steps, beta_schedule=args.ddpm_beta_schedule)
|
||||
@@ -319,7 +319,12 @@ def main(args):
|
||||
|
||||
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
|
||||
|
||||
ema_model = EMAModel(model, inv_gamma=args.ema_inv_gamma, power=args.ema_power, max_value=args.ema_max_decay)
|
||||
ema_model = EMAModel(
|
||||
accelerator.unwrap_model(model),
|
||||
inv_gamma=args.ema_inv_gamma,
|
||||
power=args.ema_power,
|
||||
max_value=args.ema_max_decay,
|
||||
)
|
||||
|
||||
# Handle the repository creation
|
||||
if accelerator.is_main_process:
|
||||
@@ -365,9 +370,9 @@ def main(args):
|
||||
# Predict the noise residual
|
||||
model_output = model(noisy_images, timesteps).sample
|
||||
|
||||
if args.predict_epsilon:
|
||||
if args.prediction_type == "epsilon":
|
||||
loss = F.mse_loss(model_output, noise) # this could have different weights!
|
||||
else:
|
||||
elif args.prediction_type == "sample":
|
||||
alpha_t = _extract_into_tensor(
|
||||
noise_scheduler.alphas_cumprod, timesteps, (clean_images.shape[0], 1, 1, 1)
|
||||
)
|
||||
@@ -376,6 +381,8 @@ def main(args):
|
||||
model_output, clean_images, reduction="none"
|
||||
) # use SNR weighting from distillation paper
|
||||
loss = loss.mean()
|
||||
else:
|
||||
raise ValueError(f"Unsupported prediction type: {args.prediction_type}")
|
||||
|
||||
accelerator.backward(loss)
|
||||
|
||||
@@ -409,11 +416,7 @@ def main(args):
|
||||
scheduler=noise_scheduler,
|
||||
)
|
||||
|
||||
deprecate("todo: remove this check", "0.10.0", "when the most used version is >= 0.8.0")
|
||||
if diffusers_version < version.parse("0.8.0"):
|
||||
generator = torch.manual_seed(0)
|
||||
else:
|
||||
generator = torch.Generator(device=pipeline.device).manual_seed(0)
|
||||
generator = torch.Generator(device=pipeline.device).manual_seed(0)
|
||||
# run pipeline in inference (sample random noise and denoise)
|
||||
images = pipeline(
|
||||
generator=generator,
|
||||
|
||||
@@ -1,6 +1,8 @@
|
||||
import argparse
|
||||
import math
|
||||
import os
|
||||
from pathlib import Path
|
||||
from typing import Optional
|
||||
|
||||
import torch
|
||||
import torch.nn.functional as F
|
||||
@@ -9,9 +11,9 @@ from accelerate import Accelerator
|
||||
from accelerate.logging import get_logger
|
||||
from datasets import load_dataset
|
||||
from diffusers import DDPMPipeline, DDPMScheduler, UNet2DModel
|
||||
from diffusers.hub_utils import init_git_repo, push_to_hub
|
||||
from diffusers.optimization import get_scheduler
|
||||
from diffusers.training_utils import EMAModel
|
||||
from huggingface_hub import HfFolder, Repository, whoami
|
||||
from onnxruntime.training.ortmodule import ORTModule
|
||||
from torchvision.transforms import (
|
||||
CenterCrop,
|
||||
@@ -28,6 +30,16 @@ from tqdm.auto import tqdm
|
||||
logger = get_logger(__name__)
|
||||
|
||||
|
||||
def get_full_repo_name(model_id: str, organization: Optional[str] = None, token: Optional[str] = None):
|
||||
if token is None:
|
||||
token = HfFolder.get_token()
|
||||
if organization is None:
|
||||
username = whoami(token)["name"]
|
||||
return f"{username}/{model_id}"
|
||||
else:
|
||||
return f"{organization}/{model_id}"
|
||||
|
||||
|
||||
def main(args):
|
||||
logging_dir = os.path.join(args.output_dir, args.logging_dir)
|
||||
accelerator = Accelerator(
|
||||
@@ -113,8 +125,22 @@ def main(args):
|
||||
|
||||
ema_model = EMAModel(model, inv_gamma=args.ema_inv_gamma, power=args.ema_power, max_value=args.ema_max_decay)
|
||||
|
||||
if args.push_to_hub:
|
||||
repo = init_git_repo(args, at_init=True)
|
||||
# Handle the repository creation
|
||||
if accelerator.is_main_process:
|
||||
if args.push_to_hub:
|
||||
if args.hub_model_id is None:
|
||||
repo_name = get_full_repo_name(Path(args.output_dir).name, token=args.hub_token)
|
||||
else:
|
||||
repo_name = args.hub_model_id
|
||||
repo = Repository(args.output_dir, clone_from=repo_name)
|
||||
|
||||
with open(os.path.join(args.output_dir, ".gitignore"), "w+") as gitignore:
|
||||
if "step_*" not in gitignore:
|
||||
gitignore.write("step_*\n")
|
||||
if "epoch_*" not in gitignore:
|
||||
gitignore.write("epoch_*\n")
|
||||
elif args.output_dir is not None:
|
||||
os.makedirs(args.output_dir, exist_ok=True)
|
||||
|
||||
if accelerator.is_main_process:
|
||||
run = os.path.split(__file__)[-1].split(".")[0]
|
||||
@@ -186,10 +212,9 @@ def main(args):
|
||||
|
||||
if epoch % args.save_model_epochs == 0 or epoch == args.num_epochs - 1:
|
||||
# save the model
|
||||
pipeline.save_pretrained(args.output_dir)
|
||||
if args.push_to_hub:
|
||||
push_to_hub(args, pipeline, repo, commit_message=f"Epoch {epoch}", blocking=False)
|
||||
else:
|
||||
pipeline.save_pretrained(args.output_dir)
|
||||
repo.push_to_hub(commit_message=f"Epoch {epoch}", blocking=False)
|
||||
accelerator.wait_for_everyone()
|
||||
|
||||
accelerator.end_training()
|
||||
|
||||
@@ -211,6 +211,7 @@ def create_unet_diffusers_config(original_config):
|
||||
"""
|
||||
Creates a config for the diffusers based on the config of the LDM model.
|
||||
"""
|
||||
model_params = original_config.model.params
|
||||
unet_params = original_config.model.params.unet_config.params
|
||||
|
||||
block_out_channels = [unet_params.model_channels * mult for mult in unet_params.channel_mult]
|
||||
@@ -230,7 +231,7 @@ def create_unet_diffusers_config(original_config):
|
||||
resolution //= 2
|
||||
|
||||
config = dict(
|
||||
sample_size=unet_params.image_size,
|
||||
sample_size=model_params.image_size,
|
||||
in_channels=unet_params.in_channels,
|
||||
out_channels=unet_params.out_channels,
|
||||
down_block_types=tuple(down_block_types),
|
||||
|
||||
@@ -215,8 +215,10 @@ def convert_models(model_path: str, output_path: str, opset: int, fp16: bool = F
|
||||
)
|
||||
del pipeline.safety_checker
|
||||
safety_checker = OnnxRuntimeModel.from_pretrained(output_path / "safety_checker")
|
||||
feature_extractor = pipeline.feature_extractor
|
||||
else:
|
||||
safety_checker = None
|
||||
feature_extractor = None
|
||||
|
||||
onnx_pipeline = OnnxStableDiffusionPipeline(
|
||||
vae_encoder=OnnxRuntimeModel.from_pretrained(output_path / "vae_encoder"),
|
||||
@@ -226,7 +228,8 @@ def convert_models(model_path: str, output_path: str, opset: int, fp16: bool = F
|
||||
unet=OnnxRuntimeModel.from_pretrained(output_path / "unet"),
|
||||
scheduler=pipeline.scheduler,
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=pipeline.feature_extractor,
|
||||
feature_extractor=feature_extractor,
|
||||
requires_safety_checker=safety_checker is not None,
|
||||
)
|
||||
|
||||
onnx_pipeline.save_pretrained(output_path)
|
||||
|
||||
6
setup.py
6
setup.py
@@ -97,6 +97,7 @@ _deps = [
|
||||
"pytest",
|
||||
"pytest-timeout",
|
||||
"pytest-xdist",
|
||||
"safetensors",
|
||||
"sentencepiece>=0.1.91,!=0.1.92",
|
||||
"scipy",
|
||||
"regex!=2019.12.17",
|
||||
@@ -184,10 +185,11 @@ extras["test"] = deps_list(
|
||||
"pytest",
|
||||
"pytest-timeout",
|
||||
"pytest-xdist",
|
||||
"safetensors",
|
||||
"sentencepiece",
|
||||
"scipy",
|
||||
"torchvision",
|
||||
"transformers"
|
||||
"transformers",
|
||||
)
|
||||
extras["torch"] = deps_list("torch", "accelerate")
|
||||
|
||||
@@ -212,7 +214,7 @@ install_requires = [
|
||||
|
||||
setup(
|
||||
name="diffusers",
|
||||
version="0.8.0", # expected format is one of x.y.z.dev0, or x.y.z.rc1 or x.y.z (no to dashes, yes to dots)
|
||||
version="0.10.0.dev0", # expected format is one of x.y.z.dev0, or x.y.z.rc1 or x.y.z (no to dashes, yes to dots)
|
||||
description="Diffusers",
|
||||
long_description=open("README.md", "r", encoding="utf-8").read(),
|
||||
long_description_content_type="text/markdown",
|
||||
|
||||
@@ -9,7 +9,7 @@ from .utils import (
|
||||
)
|
||||
|
||||
|
||||
__version__ = "0.8.0"
|
||||
__version__ = "0.10.0.dev0"
|
||||
|
||||
from .configuration_utils import ConfigMixin
|
||||
from .onnx_utils import OnnxRuntimeModel
|
||||
@@ -46,6 +46,7 @@ if is_torch_available():
|
||||
DPMSolverMultistepScheduler,
|
||||
EulerAncestralDiscreteScheduler,
|
||||
EulerDiscreteScheduler,
|
||||
HeunDiscreteScheduler,
|
||||
IPNDMScheduler,
|
||||
KarrasVeScheduler,
|
||||
PNDMScheduler,
|
||||
@@ -75,6 +76,7 @@ if is_torch_available() and is_transformers_available():
|
||||
StableDiffusionInpaintPipelineLegacy,
|
||||
StableDiffusionPipeline,
|
||||
StableDiffusionPipelineSafe,
|
||||
StableDiffusionUpscalePipeline,
|
||||
VersatileDiffusionDualGuidedPipeline,
|
||||
VersatileDiffusionImageVariationPipeline,
|
||||
VersatileDiffusionPipeline,
|
||||
|
||||
@@ -24,6 +24,8 @@ import re
|
||||
from collections import OrderedDict
|
||||
from typing import Any, Dict, Tuple, Union
|
||||
|
||||
import numpy as np
|
||||
|
||||
from huggingface_hub import hf_hub_download
|
||||
from huggingface_hub.utils import EntryNotFoundError, RepositoryNotFoundError, RevisionNotFoundError
|
||||
from requests import HTTPError
|
||||
@@ -80,20 +82,21 @@ class ConfigMixin:
|
||||
- **config_name** (`str`) -- A filename under which the config should stored when calling
|
||||
[`~ConfigMixin.save_config`] (should be overridden by parent class).
|
||||
- **ignore_for_config** (`List[str]`) -- A list of attributes that should not be saved in the config (should be
|
||||
overridden by parent class).
|
||||
- **has_compatibles** (`bool`) -- Whether the class has compatible classes (should be overridden by parent
|
||||
class).
|
||||
overridden by subclass).
|
||||
- **has_compatibles** (`bool`) -- Whether the class has compatible classes (should be overridden by subclass).
|
||||
- **_deprecated_kwargs** (`List[str]`) -- Keyword arguments that are deprecated. Note that the init function
|
||||
should only have a `kwargs` argument if at least one argument is deprecated (should be overridden by
|
||||
subclass).
|
||||
"""
|
||||
config_name = None
|
||||
ignore_for_config = []
|
||||
has_compatibles = False
|
||||
|
||||
_deprecated_kwargs = []
|
||||
|
||||
def register_to_config(self, **kwargs):
|
||||
if self.config_name is None:
|
||||
raise NotImplementedError(f"Make sure that {self.__class__} has defined a class name `config_name`")
|
||||
kwargs["_class_name"] = self.__class__.__name__
|
||||
kwargs["_diffusers_version"] = __version__
|
||||
|
||||
# Special case for `kwargs` used in deprecation warning added to schedulers
|
||||
# TODO: remove this when we remove the deprecation warning, and the `kwargs` argument,
|
||||
# or solve in a more general way.
|
||||
@@ -198,6 +201,11 @@ class ConfigMixin:
|
||||
if "dtype" in unused_kwargs:
|
||||
init_dict["dtype"] = unused_kwargs.pop("dtype")
|
||||
|
||||
# add possible deprecated kwargs
|
||||
for deprecated_kwarg in cls._deprecated_kwargs:
|
||||
if deprecated_kwarg in unused_kwargs:
|
||||
init_dict[deprecated_kwarg] = unused_kwargs.pop(deprecated_kwarg)
|
||||
|
||||
# Return model and optionally state and/or unused_kwargs
|
||||
model = cls(**init_dict)
|
||||
|
||||
@@ -462,7 +470,7 @@ class ConfigMixin:
|
||||
unused_kwargs = {**config_dict, **kwargs}
|
||||
|
||||
# 7. Define "hidden" config parameters that were saved for compatible classes
|
||||
hidden_config_dict = {k: v for k, v in original_dict.items() if k not in init_dict and not k.startswith("_")}
|
||||
hidden_config_dict = {k: v for k, v in original_dict.items() if k not in init_dict}
|
||||
|
||||
return init_dict, unused_kwargs, hidden_config_dict
|
||||
|
||||
@@ -493,6 +501,15 @@ class ConfigMixin:
|
||||
`str`: String containing all the attributes that make up this configuration instance in JSON format.
|
||||
"""
|
||||
config_dict = self._internal_dict if hasattr(self, "_internal_dict") else {}
|
||||
config_dict["_class_name"] = self.__class__.__name__
|
||||
config_dict["_diffusers_version"] = __version__
|
||||
|
||||
def to_json_saveable(value):
|
||||
if isinstance(value, np.ndarray):
|
||||
value = value.tolist()
|
||||
return value
|
||||
|
||||
config_dict = {k: to_json_saveable(v) for k, v in config_dict.items()}
|
||||
return json.dumps(config_dict, indent=2, sort_keys=True) + "\n"
|
||||
|
||||
def to_json_file(self, json_file_path: Union[str, os.PathLike]):
|
||||
@@ -520,7 +537,7 @@ def register_to_config(init):
|
||||
def inner_init(self, *args, **kwargs):
|
||||
# Ignore private kwargs in the init.
|
||||
init_kwargs = {k: v for k, v in kwargs.items() if not k.startswith("_")}
|
||||
init(self, *args, **init_kwargs)
|
||||
config_init_kwargs = {k: v for k, v in kwargs.items() if k.startswith("_")}
|
||||
if not isinstance(self, ConfigMixin):
|
||||
raise RuntimeError(
|
||||
f"`@register_for_config` was applied to {self.__class__.__name__} init method, but this class does "
|
||||
@@ -545,7 +562,9 @@ def register_to_config(init):
|
||||
if k not in ignore and k not in new_kwargs
|
||||
}
|
||||
)
|
||||
new_kwargs = {**config_init_kwargs, **new_kwargs}
|
||||
getattr(self, "register_to_config")(**new_kwargs)
|
||||
init(self, *args, **init_kwargs)
|
||||
|
||||
return inner_init
|
||||
|
||||
@@ -562,7 +581,7 @@ def flax_register_to_config(cls):
|
||||
)
|
||||
|
||||
# Ignore private kwargs in the init. Retrieve all passed attributes
|
||||
init_kwargs = {k: v for k, v in kwargs.items() if not k.startswith("_")}
|
||||
init_kwargs = {k: v for k, v in kwargs.items()}
|
||||
|
||||
# Retrieve default values
|
||||
fields = dataclasses.fields(self)
|
||||
|
||||
@@ -21,6 +21,7 @@ deps = {
|
||||
"pytest": "pytest",
|
||||
"pytest-timeout": "pytest-timeout",
|
||||
"pytest-xdist": "pytest-xdist",
|
||||
"safetensors": "safetensors",
|
||||
"sentencepiece": "sentencepiece>=0.1.91,!=0.1.92",
|
||||
"scipy": "scipy",
|
||||
"regex": "regex!=2019.12.17",
|
||||
|
||||
@@ -89,6 +89,7 @@ class ValueGuidedRLPipeline(DiffusionPipeline):
|
||||
x = x + scale * grad
|
||||
x = self.reset_x0(x, conditions, self.action_dim)
|
||||
prev_x = self.unet(x.permute(0, 2, 1), timesteps).sample.permute(0, 2, 1)
|
||||
# TODO: set prediction_type when instantiating the model
|
||||
x = self.scheduler.step(prev_x, i, x, predict_epsilon=False)["prev_sample"]
|
||||
|
||||
# apply conditions to the trajectory
|
||||
|
||||
@@ -15,16 +15,15 @@
|
||||
|
||||
|
||||
import os
|
||||
import shutil
|
||||
import sys
|
||||
from pathlib import Path
|
||||
from typing import Dict, Optional, Union
|
||||
from uuid import uuid4
|
||||
|
||||
from huggingface_hub import HfFolder, Repository, whoami
|
||||
from huggingface_hub import HfFolder, whoami
|
||||
|
||||
from . import __version__
|
||||
from .utils import ENV_VARS_TRUE_VALUES, deprecate, logging
|
||||
from .utils import ENV_VARS_TRUE_VALUES, logging
|
||||
from .utils.import_utils import (
|
||||
_flax_version,
|
||||
_jax_version,
|
||||
@@ -83,121 +82,6 @@ def get_full_repo_name(model_id: str, organization: Optional[str] = None, token:
|
||||
return f"{organization}/{model_id}"
|
||||
|
||||
|
||||
def init_git_repo(args, at_init: bool = False):
|
||||
"""
|
||||
Args:
|
||||
Initializes a git repo in `args.hub_model_id`.
|
||||
at_init (`bool`, *optional*, defaults to `False`):
|
||||
Whether this function is called before any training or not. If `self.args.overwrite_output_dir` is `True`
|
||||
and `at_init` is `True`, the path to the repo (which is `self.args.output_dir`) might be wiped out.
|
||||
"""
|
||||
deprecation_message = (
|
||||
"Please use `huggingface_hub.Repository`. "
|
||||
"See `examples/unconditional_image_generation/train_unconditional.py` for an example."
|
||||
)
|
||||
deprecate("init_git_repo()", "0.10.0", deprecation_message)
|
||||
|
||||
if hasattr(args, "local_rank") and args.local_rank not in [-1, 0]:
|
||||
return
|
||||
hub_token = args.hub_token if hasattr(args, "hub_token") else None
|
||||
use_auth_token = True if hub_token is None else hub_token
|
||||
if not hasattr(args, "hub_model_id") or args.hub_model_id is None:
|
||||
repo_name = Path(args.output_dir).absolute().name
|
||||
else:
|
||||
repo_name = args.hub_model_id
|
||||
if "/" not in repo_name:
|
||||
repo_name = get_full_repo_name(repo_name, token=hub_token)
|
||||
|
||||
try:
|
||||
repo = Repository(
|
||||
args.output_dir,
|
||||
clone_from=repo_name,
|
||||
use_auth_token=use_auth_token,
|
||||
private=args.hub_private_repo,
|
||||
)
|
||||
except EnvironmentError:
|
||||
if args.overwrite_output_dir and at_init:
|
||||
# Try again after wiping output_dir
|
||||
shutil.rmtree(args.output_dir)
|
||||
repo = Repository(
|
||||
args.output_dir,
|
||||
clone_from=repo_name,
|
||||
use_auth_token=use_auth_token,
|
||||
)
|
||||
else:
|
||||
raise
|
||||
|
||||
repo.git_pull()
|
||||
|
||||
# By default, ignore the checkpoint folders
|
||||
if not os.path.exists(os.path.join(args.output_dir, ".gitignore")):
|
||||
with open(os.path.join(args.output_dir, ".gitignore"), "w", encoding="utf-8") as writer:
|
||||
writer.writelines(["checkpoint-*/"])
|
||||
|
||||
return repo
|
||||
|
||||
|
||||
def push_to_hub(
|
||||
args,
|
||||
pipeline,
|
||||
repo: Repository,
|
||||
commit_message: Optional[str] = "End of training",
|
||||
blocking: bool = True,
|
||||
**kwargs,
|
||||
) -> str:
|
||||
"""
|
||||
Parameters:
|
||||
Upload *self.model* and *self.tokenizer* to the 🤗 model hub on the repo *self.args.hub_model_id*.
|
||||
commit_message (`str`, *optional*, defaults to `"End of training"`):
|
||||
Message to commit while pushing.
|
||||
blocking (`bool`, *optional*, defaults to `True`):
|
||||
Whether the function should return only when the `git push` has finished.
|
||||
kwargs:
|
||||
Additional keyword arguments passed along to [`create_model_card`].
|
||||
Returns:
|
||||
The url of the commit of your model in the given repository if `blocking=False`, a tuple with the url of the
|
||||
commit and an object to track the progress of the commit if `blocking=True`
|
||||
"""
|
||||
deprecation_message = (
|
||||
"Please use `huggingface_hub.Repository` and `Repository.push_to_hub()`. "
|
||||
"See `examples/unconditional_image_generation/train_unconditional.py` for an example."
|
||||
)
|
||||
deprecate("push_to_hub()", "0.10.0", deprecation_message)
|
||||
|
||||
if not hasattr(args, "hub_model_id") or args.hub_model_id is None:
|
||||
model_name = Path(args.output_dir).name
|
||||
else:
|
||||
model_name = args.hub_model_id.split("/")[-1]
|
||||
|
||||
output_dir = args.output_dir
|
||||
os.makedirs(output_dir, exist_ok=True)
|
||||
logger.info(f"Saving pipeline checkpoint to {output_dir}")
|
||||
pipeline.save_pretrained(output_dir)
|
||||
|
||||
# Only push from one node.
|
||||
if hasattr(args, "local_rank") and args.local_rank not in [-1, 0]:
|
||||
return
|
||||
|
||||
# Cancel any async push in progress if blocking=True. The commits will all be pushed together.
|
||||
if (
|
||||
blocking
|
||||
and len(repo.command_queue) > 0
|
||||
and repo.command_queue[-1] is not None
|
||||
and not repo.command_queue[-1].is_done
|
||||
):
|
||||
repo.command_queue[-1]._process.kill()
|
||||
|
||||
git_head_commit_url = repo.push_to_hub(commit_message=commit_message, blocking=blocking, auto_lfs_prune=True)
|
||||
# push separately the model card to be independent from the rest of the model
|
||||
create_model_card(args, model_name=model_name)
|
||||
try:
|
||||
repo.push_to_hub(commit_message="update model card README.md", blocking=blocking, auto_lfs_prune=True)
|
||||
except EnvironmentError as exc:
|
||||
logger.error(f"Error pushing update to the model card. Please read logs and retry.\n${exc}")
|
||||
|
||||
return git_head_commit_url
|
||||
|
||||
|
||||
def create_model_card(args, model_name):
|
||||
if not is_modelcards_available:
|
||||
raise ValueError(
|
||||
|
||||
@@ -332,7 +332,7 @@ class FlaxModelMixin:
|
||||
elif os.path.isfile(os.path.join(pretrained_path_with_subfolder, WEIGHTS_NAME)):
|
||||
raise EnvironmentError(
|
||||
f"{WEIGHTS_NAME} file found in directory {pretrained_path_with_subfolder}. Please load the model"
|
||||
" using `from_pt=True`."
|
||||
" using `from_pt=True`."
|
||||
)
|
||||
else:
|
||||
raise EnvironmentError(
|
||||
|
||||
@@ -30,8 +30,10 @@ from .utils import (
|
||||
CONFIG_NAME,
|
||||
DIFFUSERS_CACHE,
|
||||
HUGGINGFACE_CO_RESOLVE_ENDPOINT,
|
||||
SAFETENSORS_WEIGHTS_NAME,
|
||||
WEIGHTS_NAME,
|
||||
is_accelerate_available,
|
||||
is_safetensors_available,
|
||||
is_torch_version,
|
||||
logging,
|
||||
)
|
||||
@@ -51,6 +53,9 @@ if is_accelerate_available():
|
||||
from accelerate.utils import set_module_tensor_to_device
|
||||
from accelerate.utils.versions import is_torch_version
|
||||
|
||||
if is_safetensors_available():
|
||||
import safetensors
|
||||
|
||||
|
||||
def get_parameter_device(parameter: torch.nn.Module):
|
||||
try:
|
||||
@@ -84,10 +89,13 @@ def get_parameter_dtype(parameter: torch.nn.Module):
|
||||
|
||||
def load_state_dict(checkpoint_file: Union[str, os.PathLike]):
|
||||
"""
|
||||
Reads a PyTorch checkpoint file, returning properly formatted errors if they arise.
|
||||
Reads a checkpoint file, returning properly formatted errors if they arise.
|
||||
"""
|
||||
try:
|
||||
return torch.load(checkpoint_file, map_location="cpu")
|
||||
if os.path.basename(checkpoint_file) == WEIGHTS_NAME:
|
||||
return torch.load(checkpoint_file, map_location="cpu")
|
||||
else:
|
||||
return safetensors.torch.load_file(checkpoint_file, device="cpu")
|
||||
except Exception as e:
|
||||
try:
|
||||
with open(checkpoint_file) as f:
|
||||
@@ -104,7 +112,7 @@ def load_state_dict(checkpoint_file: Union[str, os.PathLike]):
|
||||
) from e
|
||||
except (UnicodeDecodeError, ValueError):
|
||||
raise OSError(
|
||||
f"Unable to load weights from pytorch checkpoint file for '{checkpoint_file}' "
|
||||
f"Unable to load weights from checkpoint file for '{checkpoint_file}' "
|
||||
f"at '{checkpoint_file}'. "
|
||||
"If you tried to load a PyTorch model from a TF 2.0 checkpoint, please set from_tf=True."
|
||||
)
|
||||
@@ -375,80 +383,44 @@ class ModelMixin(torch.nn.Module):
|
||||
|
||||
# This variable will flag if we're loading a sharded checkpoint. In this case the archive file is just the
|
||||
# Load model
|
||||
pretrained_model_name_or_path = str(pretrained_model_name_or_path)
|
||||
if os.path.isdir(pretrained_model_name_or_path):
|
||||
if os.path.isfile(os.path.join(pretrained_model_name_or_path, WEIGHTS_NAME)):
|
||||
# Load from a PyTorch checkpoint
|
||||
model_file = os.path.join(pretrained_model_name_or_path, WEIGHTS_NAME)
|
||||
elif subfolder is not None and os.path.isfile(
|
||||
os.path.join(pretrained_model_name_or_path, subfolder, WEIGHTS_NAME)
|
||||
):
|
||||
model_file = os.path.join(pretrained_model_name_or_path, subfolder, WEIGHTS_NAME)
|
||||
else:
|
||||
raise EnvironmentError(
|
||||
f"Error no file named {WEIGHTS_NAME} found in directory {pretrained_model_name_or_path}."
|
||||
)
|
||||
else:
|
||||
|
||||
model_file = None
|
||||
if is_safetensors_available():
|
||||
try:
|
||||
# Load from URL or cache if already cached
|
||||
model_file = hf_hub_download(
|
||||
model_file = _get_model_file(
|
||||
pretrained_model_name_or_path,
|
||||
filename=WEIGHTS_NAME,
|
||||
weights_name=SAFETENSORS_WEIGHTS_NAME,
|
||||
cache_dir=cache_dir,
|
||||
force_download=force_download,
|
||||
proxies=proxies,
|
||||
resume_download=resume_download,
|
||||
proxies=proxies,
|
||||
local_files_only=local_files_only,
|
||||
use_auth_token=use_auth_token,
|
||||
user_agent=user_agent,
|
||||
subfolder=subfolder,
|
||||
revision=revision,
|
||||
subfolder=subfolder,
|
||||
user_agent=user_agent,
|
||||
)
|
||||
|
||||
except RepositoryNotFoundError:
|
||||
raise EnvironmentError(
|
||||
f"{pretrained_model_name_or_path} is not a local folder and is not a valid model identifier "
|
||||
"listed on 'https://huggingface.co/models'\nIf this is a private repository, make sure to pass a "
|
||||
"token having permission to this repo with `use_auth_token` or log in with `huggingface-cli "
|
||||
"login`."
|
||||
)
|
||||
except RevisionNotFoundError:
|
||||
raise EnvironmentError(
|
||||
f"{revision} is not a valid git identifier (branch name, tag name or commit id) that exists for "
|
||||
"this model name. Check the model page at "
|
||||
f"'https://huggingface.co/{pretrained_model_name_or_path}' for available revisions."
|
||||
)
|
||||
except EntryNotFoundError:
|
||||
raise EnvironmentError(
|
||||
f"{pretrained_model_name_or_path} does not appear to have a file named {WEIGHTS_NAME}."
|
||||
)
|
||||
except HTTPError as err:
|
||||
raise EnvironmentError(
|
||||
"There was a specific connection error when trying to load"
|
||||
f" {pretrained_model_name_or_path}:\n{err}"
|
||||
)
|
||||
except ValueError:
|
||||
raise EnvironmentError(
|
||||
f"We couldn't connect to '{HUGGINGFACE_CO_RESOLVE_ENDPOINT}' to load this model, couldn't find it"
|
||||
f" in the cached files and it looks like {pretrained_model_name_or_path} is not the path to a"
|
||||
f" directory containing a file named {WEIGHTS_NAME} or"
|
||||
" \nCheckout your internet connection or see how to run the library in"
|
||||
" offline mode at 'https://huggingface.co/docs/diffusers/installation#offline-mode'."
|
||||
)
|
||||
except EnvironmentError:
|
||||
raise EnvironmentError(
|
||||
f"Can't load the model for '{pretrained_model_name_or_path}'. If you were trying to load it from "
|
||||
"'https://huggingface.co/models', make sure you don't have a local directory with the same name. "
|
||||
f"Otherwise, make sure '{pretrained_model_name_or_path}' is the correct path to a directory "
|
||||
f"containing a file named {WEIGHTS_NAME}"
|
||||
)
|
||||
|
||||
# restore default dtype
|
||||
except:
|
||||
pass
|
||||
if model_file is None:
|
||||
model_file = _get_model_file(
|
||||
pretrained_model_name_or_path,
|
||||
weights_name=WEIGHTS_NAME,
|
||||
cache_dir=cache_dir,
|
||||
force_download=force_download,
|
||||
resume_download=resume_download,
|
||||
proxies=proxies,
|
||||
local_files_only=local_files_only,
|
||||
use_auth_token=use_auth_token,
|
||||
revision=revision,
|
||||
subfolder=subfolder,
|
||||
user_agent=user_agent,
|
||||
)
|
||||
|
||||
if low_cpu_mem_usage:
|
||||
# Instantiate model with empty weights
|
||||
with accelerate.init_empty_weights():
|
||||
model, unused_kwargs = cls.from_config(
|
||||
config, unused_kwargs = cls.load_config(
|
||||
config_path,
|
||||
cache_dir=cache_dir,
|
||||
return_unused_kwargs=True,
|
||||
@@ -462,6 +434,7 @@ class ModelMixin(torch.nn.Module):
|
||||
device_map=device_map,
|
||||
**kwargs,
|
||||
)
|
||||
model = cls.from_config(config, **unused_kwargs)
|
||||
|
||||
# if device_map is Non,e load the state dict on move the params from meta device to the cpu
|
||||
if device_map is None:
|
||||
@@ -482,7 +455,7 @@ class ModelMixin(torch.nn.Module):
|
||||
"error_msgs": [],
|
||||
}
|
||||
else:
|
||||
model, unused_kwargs = cls.from_config(
|
||||
config, unused_kwargs = cls.load_config(
|
||||
config_path,
|
||||
cache_dir=cache_dir,
|
||||
return_unused_kwargs=True,
|
||||
@@ -496,8 +469,24 @@ class ModelMixin(torch.nn.Module):
|
||||
device_map=device_map,
|
||||
**kwargs,
|
||||
)
|
||||
model = cls.from_config(config, **unused_kwargs)
|
||||
|
||||
state_dict = load_state_dict(model_file)
|
||||
dtype = set(v.dtype for v in state_dict.values())
|
||||
|
||||
if len(dtype) > 1 and torch.float32 not in dtype:
|
||||
raise ValueError(
|
||||
f"The weights of the model file {model_file} have a mixture of incompatible dtypes {dtype}. Please"
|
||||
f" make sure that {model_file} weights have only one dtype."
|
||||
)
|
||||
elif len(dtype) > 1 and torch.float32 in dtype:
|
||||
dtype = torch.float32
|
||||
else:
|
||||
dtype = dtype.pop()
|
||||
|
||||
# move model to correct dtype
|
||||
model = model.to(dtype)
|
||||
|
||||
model, missing_keys, unexpected_keys, mismatched_keys, error_msgs = cls._load_pretrained_model(
|
||||
model,
|
||||
state_dict,
|
||||
@@ -677,15 +666,86 @@ class ModelMixin(torch.nn.Module):
|
||||
return sum(p.numel() for p in self.parameters() if p.requires_grad or not only_trainable)
|
||||
|
||||
|
||||
def unwrap_model(model: torch.nn.Module) -> torch.nn.Module:
|
||||
"""
|
||||
Recursively unwraps a model from potential containers (as used in distributed training).
|
||||
|
||||
Args:
|
||||
model (`torch.nn.Module`): The model to unwrap.
|
||||
"""
|
||||
# since there could be multiple levels of wrapping, unwrap recursively
|
||||
if hasattr(model, "module"):
|
||||
return unwrap_model(model.module)
|
||||
def _get_model_file(
|
||||
pretrained_model_name_or_path,
|
||||
*,
|
||||
weights_name,
|
||||
subfolder,
|
||||
cache_dir,
|
||||
force_download,
|
||||
proxies,
|
||||
resume_download,
|
||||
local_files_only,
|
||||
use_auth_token,
|
||||
user_agent,
|
||||
revision,
|
||||
):
|
||||
pretrained_model_name_or_path = str(pretrained_model_name_or_path)
|
||||
if os.path.isdir(pretrained_model_name_or_path):
|
||||
if os.path.isfile(os.path.join(pretrained_model_name_or_path, weights_name)):
|
||||
# Load from a PyTorch checkpoint
|
||||
model_file = os.path.join(pretrained_model_name_or_path, weights_name)
|
||||
return model_file
|
||||
elif subfolder is not None and os.path.isfile(
|
||||
os.path.join(pretrained_model_name_or_path, subfolder, weights_name)
|
||||
):
|
||||
model_file = os.path.join(pretrained_model_name_or_path, subfolder, weights_name)
|
||||
return model_file
|
||||
else:
|
||||
raise EnvironmentError(
|
||||
f"Error no file named {weights_name} found in directory {pretrained_model_name_or_path}."
|
||||
)
|
||||
else:
|
||||
return model
|
||||
try:
|
||||
# Load from URL or cache if already cached
|
||||
model_file = hf_hub_download(
|
||||
pretrained_model_name_or_path,
|
||||
filename=weights_name,
|
||||
cache_dir=cache_dir,
|
||||
force_download=force_download,
|
||||
proxies=proxies,
|
||||
resume_download=resume_download,
|
||||
local_files_only=local_files_only,
|
||||
use_auth_token=use_auth_token,
|
||||
user_agent=user_agent,
|
||||
subfolder=subfolder,
|
||||
revision=revision,
|
||||
)
|
||||
return model_file
|
||||
|
||||
except RepositoryNotFoundError:
|
||||
raise EnvironmentError(
|
||||
f"{pretrained_model_name_or_path} is not a local folder and is not a valid model identifier "
|
||||
"listed on 'https://huggingface.co/models'\nIf this is a private repository, make sure to pass a "
|
||||
"token having permission to this repo with `use_auth_token` or log in with `huggingface-cli "
|
||||
"login`."
|
||||
)
|
||||
except RevisionNotFoundError:
|
||||
raise EnvironmentError(
|
||||
f"{revision} is not a valid git identifier (branch name, tag name or commit id) that exists for "
|
||||
"this model name. Check the model page at "
|
||||
f"'https://huggingface.co/{pretrained_model_name_or_path}' for available revisions."
|
||||
)
|
||||
except EntryNotFoundError:
|
||||
raise EnvironmentError(
|
||||
f"{pretrained_model_name_or_path} does not appear to have a file named {weights_name}."
|
||||
)
|
||||
except HTTPError as err:
|
||||
raise EnvironmentError(
|
||||
f"There was a specific connection error when trying to load {pretrained_model_name_or_path}:\n{err}"
|
||||
)
|
||||
except ValueError:
|
||||
raise EnvironmentError(
|
||||
f"We couldn't connect to '{HUGGINGFACE_CO_RESOLVE_ENDPOINT}' to load this model, couldn't find it"
|
||||
f" in the cached files and it looks like {pretrained_model_name_or_path} is not the path to a"
|
||||
f" directory containing a file named {weights_name} or"
|
||||
" \nCheckout your internet connection or see how to run the library in"
|
||||
" offline mode at 'https://huggingface.co/docs/diffusers/installation#offline-mode'."
|
||||
)
|
||||
except EnvironmentError:
|
||||
raise EnvironmentError(
|
||||
f"Can't load the model for '{pretrained_model_name_or_path}'. If you were trying to load it from "
|
||||
"'https://huggingface.co/models', make sure you don't have a local directory with the same name. "
|
||||
f"Otherwise, make sure '{pretrained_model_name_or_path}' is the correct path to a directory "
|
||||
f"containing a file named {weights_name}"
|
||||
)
|
||||
|
||||
@@ -100,6 +100,7 @@ class Transformer2DModel(ModelMixin, ConfigMixin):
|
||||
activation_fn: str = "geglu",
|
||||
num_embeds_ada_norm: Optional[int] = None,
|
||||
use_linear_projection: bool = False,
|
||||
only_cross_attention: bool = False,
|
||||
):
|
||||
super().__init__()
|
||||
self.use_linear_projection = use_linear_projection
|
||||
@@ -157,6 +158,7 @@ class Transformer2DModel(ModelMixin, ConfigMixin):
|
||||
activation_fn=activation_fn,
|
||||
num_embeds_ada_norm=num_embeds_ada_norm,
|
||||
attention_bias=attention_bias,
|
||||
only_cross_attention=only_cross_attention,
|
||||
)
|
||||
for d in range(num_layers)
|
||||
]
|
||||
@@ -219,11 +221,15 @@ class Transformer2DModel(ModelMixin, ConfigMixin):
|
||||
# 3. Output
|
||||
if self.is_input_continuous:
|
||||
if not self.use_linear_projection:
|
||||
hidden_states = hidden_states.reshape(batch, height, weight, inner_dim).permute(0, 3, 1, 2)
|
||||
hidden_states = (
|
||||
hidden_states.reshape(batch, height, weight, inner_dim).permute(0, 3, 1, 2).contiguous()
|
||||
)
|
||||
hidden_states = self.proj_out(hidden_states)
|
||||
else:
|
||||
hidden_states = self.proj_out(hidden_states)
|
||||
hidden_states = hidden_states.reshape(batch, height, weight, inner_dim).permute(0, 3, 1, 2)
|
||||
hidden_states = (
|
||||
hidden_states.reshape(batch, height, weight, inner_dim).permute(0, 3, 1, 2).contiguous()
|
||||
)
|
||||
|
||||
output = hidden_states + residual
|
||||
elif self.is_input_vectorized:
|
||||
@@ -284,11 +290,19 @@ class AttentionBlock(nn.Module):
|
||||
self.rescale_output_factor = rescale_output_factor
|
||||
self.proj_attn = nn.Linear(channels, channels, 1)
|
||||
|
||||
def transpose_for_scores(self, projection: torch.Tensor) -> torch.Tensor:
|
||||
new_projection_shape = projection.size()[:-1] + (self.num_heads, -1)
|
||||
# move heads to 2nd position (B, T, H * D) -> (B, T, H, D) -> (B, H, T, D)
|
||||
new_projection = projection.view(new_projection_shape).permute(0, 2, 1, 3)
|
||||
return new_projection
|
||||
def reshape_heads_to_batch_dim(self, tensor):
|
||||
batch_size, seq_len, dim = tensor.shape
|
||||
head_size = self.num_heads
|
||||
tensor = tensor.reshape(batch_size, seq_len, head_size, dim // head_size)
|
||||
tensor = tensor.permute(0, 2, 1, 3).reshape(batch_size * head_size, seq_len, dim // head_size)
|
||||
return tensor
|
||||
|
||||
def reshape_batch_dim_to_heads(self, tensor):
|
||||
batch_size, seq_len, dim = tensor.shape
|
||||
head_size = self.num_heads
|
||||
tensor = tensor.reshape(batch_size // head_size, head_size, seq_len, dim)
|
||||
tensor = tensor.permute(0, 2, 1, 3).reshape(batch_size // head_size, seq_len, dim * head_size)
|
||||
return tensor
|
||||
|
||||
def forward(self, hidden_states):
|
||||
residual = hidden_states
|
||||
@@ -306,50 +320,28 @@ class AttentionBlock(nn.Module):
|
||||
|
||||
scale = 1 / math.sqrt(self.channels / self.num_heads)
|
||||
|
||||
# get scores
|
||||
if self.num_heads > 1:
|
||||
query_states = self.transpose_for_scores(query_proj)
|
||||
key_states = self.transpose_for_scores(key_proj)
|
||||
value_states = self.transpose_for_scores(value_proj)
|
||||
|
||||
# TODO: is there a way to perform batched matmul (e.g. baddbmm) on 4D tensors?
|
||||
# or reformulate this into a 3D problem?
|
||||
# TODO: measure whether on MPS device it would be faster to do this matmul via einsum
|
||||
# as some matmuls can be 1.94x slower than an equivalent einsum on MPS
|
||||
# https://gist.github.com/Birch-san/cba16789ec27bb20996a4b4831b13ce0
|
||||
attention_scores = torch.matmul(query_states, key_states.transpose(-1, -2)) * scale
|
||||
else:
|
||||
query_states, key_states, value_states = query_proj, key_proj, value_proj
|
||||
|
||||
attention_scores = torch.baddbmm(
|
||||
torch.empty(
|
||||
query_states.shape[0],
|
||||
query_states.shape[1],
|
||||
key_states.shape[1],
|
||||
dtype=query_states.dtype,
|
||||
device=query_states.device,
|
||||
),
|
||||
query_states,
|
||||
key_states.transpose(-1, -2),
|
||||
beta=0,
|
||||
alpha=scale,
|
||||
)
|
||||
query_proj = self.reshape_heads_to_batch_dim(query_proj)
|
||||
key_proj = self.reshape_heads_to_batch_dim(key_proj)
|
||||
value_proj = self.reshape_heads_to_batch_dim(value_proj)
|
||||
|
||||
attention_scores = torch.baddbmm(
|
||||
torch.empty(
|
||||
query_proj.shape[0],
|
||||
query_proj.shape[1],
|
||||
key_proj.shape[1],
|
||||
dtype=query_proj.dtype,
|
||||
device=query_proj.device,
|
||||
),
|
||||
query_proj,
|
||||
key_proj.transpose(-1, -2),
|
||||
beta=0,
|
||||
alpha=scale,
|
||||
)
|
||||
attention_probs = torch.softmax(attention_scores.float(), dim=-1).type(attention_scores.dtype)
|
||||
hidden_states = torch.bmm(attention_probs, value_proj)
|
||||
|
||||
# compute attention output
|
||||
if self.num_heads > 1:
|
||||
# TODO: is there a way to perform batched matmul (e.g. bmm) on 4D tensors?
|
||||
# or reformulate this into a 3D problem?
|
||||
# TODO: measure whether on MPS device it would be faster to do this matmul via einsum
|
||||
# as some matmuls can be 1.94x slower than an equivalent einsum on MPS
|
||||
# https://gist.github.com/Birch-san/cba16789ec27bb20996a4b4831b13ce0
|
||||
hidden_states = torch.matmul(attention_probs, value_states)
|
||||
hidden_states = hidden_states.permute(0, 2, 1, 3).contiguous()
|
||||
new_hidden_states_shape = hidden_states.size()[:-2] + (self.channels,)
|
||||
hidden_states = hidden_states.view(new_hidden_states_shape)
|
||||
else:
|
||||
hidden_states = torch.bmm(attention_probs, value_states)
|
||||
# reshape hidden_states
|
||||
hidden_states = self.reshape_batch_dim_to_heads(hidden_states)
|
||||
|
||||
# compute next hidden_states
|
||||
hidden_states = self.proj_attn(hidden_states)
|
||||
@@ -387,14 +379,17 @@ class BasicTransformerBlock(nn.Module):
|
||||
activation_fn: str = "geglu",
|
||||
num_embeds_ada_norm: Optional[int] = None,
|
||||
attention_bias: bool = False,
|
||||
only_cross_attention: bool = False,
|
||||
):
|
||||
super().__init__()
|
||||
self.only_cross_attention = only_cross_attention
|
||||
self.attn1 = CrossAttention(
|
||||
query_dim=dim,
|
||||
heads=num_attention_heads,
|
||||
dim_head=attention_head_dim,
|
||||
dropout=dropout,
|
||||
bias=attention_bias,
|
||||
cross_attention_dim=cross_attention_dim if only_cross_attention else None,
|
||||
) # is a self-attention
|
||||
self.ff = FeedForward(dim, dropout=dropout, activation_fn=activation_fn)
|
||||
self.attn2 = CrossAttention(
|
||||
@@ -461,7 +456,11 @@ class BasicTransformerBlock(nn.Module):
|
||||
norm_hidden_states = (
|
||||
self.norm1(hidden_states, timestep) if self.use_ada_layer_norm else self.norm1(hidden_states)
|
||||
)
|
||||
hidden_states = self.attn1(norm_hidden_states) + hidden_states
|
||||
|
||||
if self.only_cross_attention:
|
||||
hidden_states = self.attn1(norm_hidden_states, context) + hidden_states
|
||||
else:
|
||||
hidden_states = self.attn1(norm_hidden_states) + hidden_states
|
||||
|
||||
# 2. Cross-Attention
|
||||
norm_hidden_states = (
|
||||
|
||||
@@ -104,6 +104,8 @@ class FlaxBasicTransformerBlock(nn.Module):
|
||||
Hidden states dimension inside each head
|
||||
dropout (:obj:`float`, *optional*, defaults to 0.0):
|
||||
Dropout rate
|
||||
only_cross_attention (`bool`, defaults to `False`):
|
||||
Whether to only apply cross attention.
|
||||
dtype (:obj:`jnp.dtype`, *optional*, defaults to jnp.float32):
|
||||
Parameters `dtype`
|
||||
"""
|
||||
@@ -111,10 +113,11 @@ class FlaxBasicTransformerBlock(nn.Module):
|
||||
n_heads: int
|
||||
d_head: int
|
||||
dropout: float = 0.0
|
||||
only_cross_attention: bool = False
|
||||
dtype: jnp.dtype = jnp.float32
|
||||
|
||||
def setup(self):
|
||||
# self attention
|
||||
# self attention (or cross_attention if only_cross_attention is True)
|
||||
self.attn1 = FlaxAttentionBlock(self.dim, self.n_heads, self.d_head, self.dropout, dtype=self.dtype)
|
||||
# cross attention
|
||||
self.attn2 = FlaxAttentionBlock(self.dim, self.n_heads, self.d_head, self.dropout, dtype=self.dtype)
|
||||
@@ -126,7 +129,10 @@ class FlaxBasicTransformerBlock(nn.Module):
|
||||
def __call__(self, hidden_states, context, deterministic=True):
|
||||
# self attention
|
||||
residual = hidden_states
|
||||
hidden_states = self.attn1(self.norm1(hidden_states), deterministic=deterministic)
|
||||
if self.only_cross_attention:
|
||||
hidden_states = self.attn1(self.norm1(hidden_states), context, deterministic=deterministic)
|
||||
else:
|
||||
hidden_states = self.attn1(self.norm1(hidden_states), deterministic=deterministic)
|
||||
hidden_states = hidden_states + residual
|
||||
|
||||
# cross attention
|
||||
@@ -159,6 +165,8 @@ class FlaxTransformer2DModel(nn.Module):
|
||||
Number of transformers block
|
||||
dropout (:obj:`float`, *optional*, defaults to 0.0):
|
||||
Dropout rate
|
||||
use_linear_projection (`bool`, defaults to `False`): tbd
|
||||
only_cross_attention (`bool`, defaults to `False`): tbd
|
||||
dtype (:obj:`jnp.dtype`, *optional*, defaults to jnp.float32):
|
||||
Parameters `dtype`
|
||||
"""
|
||||
@@ -167,49 +175,70 @@ class FlaxTransformer2DModel(nn.Module):
|
||||
d_head: int
|
||||
depth: int = 1
|
||||
dropout: float = 0.0
|
||||
use_linear_projection: bool = False
|
||||
only_cross_attention: bool = False
|
||||
dtype: jnp.dtype = jnp.float32
|
||||
|
||||
def setup(self):
|
||||
self.norm = nn.GroupNorm(num_groups=32, epsilon=1e-5)
|
||||
|
||||
inner_dim = self.n_heads * self.d_head
|
||||
self.proj_in = nn.Conv(
|
||||
inner_dim,
|
||||
kernel_size=(1, 1),
|
||||
strides=(1, 1),
|
||||
padding="VALID",
|
||||
dtype=self.dtype,
|
||||
)
|
||||
if self.use_linear_projection:
|
||||
self.proj_in = nn.Dense(inner_dim, dtype=self.dtype)
|
||||
else:
|
||||
self.proj_in = nn.Conv(
|
||||
inner_dim,
|
||||
kernel_size=(1, 1),
|
||||
strides=(1, 1),
|
||||
padding="VALID",
|
||||
dtype=self.dtype,
|
||||
)
|
||||
|
||||
self.transformer_blocks = [
|
||||
FlaxBasicTransformerBlock(inner_dim, self.n_heads, self.d_head, dropout=self.dropout, dtype=self.dtype)
|
||||
FlaxBasicTransformerBlock(
|
||||
inner_dim,
|
||||
self.n_heads,
|
||||
self.d_head,
|
||||
dropout=self.dropout,
|
||||
only_cross_attention=self.only_cross_attention,
|
||||
dtype=self.dtype,
|
||||
)
|
||||
for _ in range(self.depth)
|
||||
]
|
||||
|
||||
self.proj_out = nn.Conv(
|
||||
inner_dim,
|
||||
kernel_size=(1, 1),
|
||||
strides=(1, 1),
|
||||
padding="VALID",
|
||||
dtype=self.dtype,
|
||||
)
|
||||
if self.use_linear_projection:
|
||||
self.proj_out = nn.Dense(inner_dim, dtype=self.dtype)
|
||||
else:
|
||||
self.proj_out = nn.Conv(
|
||||
inner_dim,
|
||||
kernel_size=(1, 1),
|
||||
strides=(1, 1),
|
||||
padding="VALID",
|
||||
dtype=self.dtype,
|
||||
)
|
||||
|
||||
def __call__(self, hidden_states, context, deterministic=True):
|
||||
batch, height, width, channels = hidden_states.shape
|
||||
residual = hidden_states
|
||||
hidden_states = self.norm(hidden_states)
|
||||
hidden_states = self.proj_in(hidden_states)
|
||||
|
||||
hidden_states = hidden_states.reshape(batch, height * width, channels)
|
||||
if self.use_linear_projection:
|
||||
hidden_states = hidden_states.reshape(batch, height * width, channels)
|
||||
hidden_states = self.proj_in(hidden_states)
|
||||
else:
|
||||
hidden_states = self.proj_in(hidden_states)
|
||||
hidden_states = hidden_states.reshape(batch, height * width, channels)
|
||||
|
||||
for transformer_block in self.transformer_blocks:
|
||||
hidden_states = transformer_block(hidden_states, context, deterministic=deterministic)
|
||||
|
||||
hidden_states = hidden_states.reshape(batch, height, width, channels)
|
||||
if self.use_linear_projection:
|
||||
hidden_states = self.proj_out(hidden_states)
|
||||
hidden_states = hidden_states.reshape(batch, height, width, channels)
|
||||
else:
|
||||
hidden_states = hidden_states.reshape(batch, height, width, channels)
|
||||
hidden_states = self.proj_out(hidden_states)
|
||||
|
||||
hidden_states = self.proj_out(hidden_states)
|
||||
hidden_states = hidden_states + residual
|
||||
|
||||
return hidden_states
|
||||
|
||||
|
||||
|
||||
@@ -288,8 +288,16 @@ _kernels = {
|
||||
}
|
||||
|
||||
|
||||
class Downsample1d(nn.Module):
|
||||
def __init__(self, kernel="linear", pad_mode="reflect"):
|
||||
class KernelDownsample1D(nn.Module):
|
||||
"""
|
||||
A static downsample module that is not updated by the optimizer.
|
||||
|
||||
Parameters:
|
||||
kernel (`str`): `linear`, `cubic`, or `lanczos3` for different static kernels used in convolution.
|
||||
pad_mode (`str`): defaults to `reflect`, use with torch.nn.functional.pad.
|
||||
"""
|
||||
|
||||
def __init__(self, kernel: str = "linear", pad_mode: str = "reflect"):
|
||||
super().__init__()
|
||||
self.pad_mode = pad_mode
|
||||
kernel_1d = torch.tensor(_kernels[kernel])
|
||||
@@ -304,8 +312,16 @@ class Downsample1d(nn.Module):
|
||||
return F.conv1d(hidden_states, weight, stride=2)
|
||||
|
||||
|
||||
class Upsample1d(nn.Module):
|
||||
def __init__(self, kernel="linear", pad_mode="reflect"):
|
||||
class KernelUpsample1D(nn.Module):
|
||||
"""
|
||||
A static upsample module that is not updated by the optimizer.
|
||||
|
||||
Parameters:
|
||||
kernel (`str`): `linear`, `cubic`, or `lanczos3` for different static kernels used in convolution.
|
||||
pad_mode (`str`): defaults to `reflect`, use with torch.nn.functional.pad.
|
||||
"""
|
||||
|
||||
def __init__(self, kernel: str = "linear", pad_mode: str = "reflect"):
|
||||
super().__init__()
|
||||
self.pad_mode = pad_mode
|
||||
kernel_1d = torch.tensor(_kernels[kernel]) * 2
|
||||
@@ -321,7 +337,7 @@ class Upsample1d(nn.Module):
|
||||
|
||||
|
||||
class SelfAttention1d(nn.Module):
|
||||
def __init__(self, in_channels, n_head=1, dropout_rate=0.0):
|
||||
def __init__(self, in_channels: int, n_head: int = 1, dropout_rate: float = 0.0):
|
||||
super().__init__()
|
||||
self.channels = in_channels
|
||||
self.group_norm = nn.GroupNorm(1, num_channels=in_channels)
|
||||
@@ -379,7 +395,7 @@ class SelfAttention1d(nn.Module):
|
||||
|
||||
|
||||
class ResConvBlock(nn.Module):
|
||||
def __init__(self, in_channels, mid_channels, out_channels, is_last=False):
|
||||
def __init__(self, in_channels: int, mid_channels: int, out_channels: int, is_last: bool = False):
|
||||
super().__init__()
|
||||
self.is_last = is_last
|
||||
self.has_conv_skip = in_channels != out_channels
|
||||
@@ -413,13 +429,12 @@ class ResConvBlock(nn.Module):
|
||||
|
||||
|
||||
class UNetMidBlock1D(nn.Module):
|
||||
def __init__(self, mid_channels, in_channels, out_channels=None):
|
||||
def __init__(self, mid_channels: int, in_channels: int, out_channels: int = None):
|
||||
super().__init__()
|
||||
|
||||
out_channels = in_channels if out_channels is None else out_channels
|
||||
|
||||
# there is always at least one resnet
|
||||
self.down = Downsample1d("cubic")
|
||||
self.down = KernelDownsample1D("cubic")
|
||||
resnets = [
|
||||
ResConvBlock(in_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, mid_channels),
|
||||
@@ -436,7 +451,7 @@ class UNetMidBlock1D(nn.Module):
|
||||
SelfAttention1d(mid_channels, mid_channels // 32),
|
||||
SelfAttention1d(out_channels, out_channels // 32),
|
||||
]
|
||||
self.up = Upsample1d(kernel="cubic")
|
||||
self.up = KernelUpsample1D(kernel="cubic")
|
||||
|
||||
self.attentions = nn.ModuleList(attentions)
|
||||
self.resnets = nn.ModuleList(resnets)
|
||||
@@ -453,21 +468,26 @@ class UNetMidBlock1D(nn.Module):
|
||||
|
||||
|
||||
class AttnDownBlock1D(nn.Module):
|
||||
def __init__(self, out_channels, in_channels, mid_channels=None):
|
||||
def __init__(self, out_channels: int, in_channels: int, num_layers: int = 3, mid_channels: int = None):
|
||||
super().__init__()
|
||||
|
||||
if num_layers < 1:
|
||||
raise ValueError("AttnDownBlock1D requires added num_layers >= 1")
|
||||
|
||||
mid_channels = out_channels if mid_channels is None else mid_channels
|
||||
|
||||
self.down = Downsample1d("cubic")
|
||||
resnets = [
|
||||
ResConvBlock(in_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, out_channels),
|
||||
]
|
||||
attentions = [
|
||||
SelfAttention1d(mid_channels, mid_channels // 32),
|
||||
SelfAttention1d(mid_channels, mid_channels // 32),
|
||||
SelfAttention1d(out_channels, out_channels // 32),
|
||||
]
|
||||
self.down = KernelDownsample1D("cubic")
|
||||
resnets = []
|
||||
attentions = []
|
||||
|
||||
for i in range(num_layers):
|
||||
in_channels = in_channels if i == 0 else mid_channels
|
||||
if i < (num_layers - 1):
|
||||
resnets.append(ResConvBlock(in_channels, mid_channels, mid_channels))
|
||||
attentions.append(SelfAttention1d(mid_channels, mid_channels // 32))
|
||||
else:
|
||||
resnets.append(ResConvBlock(mid_channels, mid_channels, out_channels))
|
||||
attentions.append(SelfAttention1d(out_channels, out_channels // 32))
|
||||
|
||||
self.attentions = nn.ModuleList(attentions)
|
||||
self.resnets = nn.ModuleList(resnets)
|
||||
@@ -483,16 +503,22 @@ class AttnDownBlock1D(nn.Module):
|
||||
|
||||
|
||||
class DownBlock1D(nn.Module):
|
||||
def __init__(self, out_channels, in_channels, mid_channels=None):
|
||||
def __init__(self, out_channels: int, in_channels: int, mid_channels: int = None, num_layers: int = 3):
|
||||
super().__init__()
|
||||
if num_layers < 1:
|
||||
raise ValueError("DownBlock1D requires added num_layers >= 1")
|
||||
|
||||
mid_channels = out_channels if mid_channels is None else mid_channels
|
||||
|
||||
self.down = Downsample1d("cubic")
|
||||
resnets = [
|
||||
ResConvBlock(in_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, out_channels),
|
||||
]
|
||||
self.down = KernelDownsample1D("cubic")
|
||||
resnets = []
|
||||
|
||||
for i in range(num_layers):
|
||||
in_channels = in_channels if i == 0 else mid_channels
|
||||
if i < (num_layers - 1):
|
||||
resnets.append(ResConvBlock(in_channels, mid_channels, mid_channels))
|
||||
else:
|
||||
resnets.append(ResConvBlock(mid_channels, mid_channels, out_channels))
|
||||
|
||||
self.resnets = nn.ModuleList(resnets)
|
||||
|
||||
@@ -506,15 +532,21 @@ class DownBlock1D(nn.Module):
|
||||
|
||||
|
||||
class DownBlock1DNoSkip(nn.Module):
|
||||
def __init__(self, out_channels, in_channels, mid_channels=None):
|
||||
def __init__(self, out_channels: int, in_channels: int, mid_channels: int = None, num_layers: int = 3):
|
||||
super().__init__()
|
||||
if num_layers < 1:
|
||||
raise ValueError("DownBlock1DNoSkip requires added num_layers >= 1")
|
||||
|
||||
mid_channels = out_channels if mid_channels is None else mid_channels
|
||||
|
||||
resnets = [
|
||||
ResConvBlock(in_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, out_channels),
|
||||
]
|
||||
resnets = []
|
||||
|
||||
for i in range(num_layers):
|
||||
in_channels = in_channels if i == 0 else mid_channels
|
||||
if i < (num_layers - 1):
|
||||
resnets.append(ResConvBlock(in_channels, mid_channels, mid_channels))
|
||||
else:
|
||||
resnets.append(ResConvBlock(mid_channels, mid_channels, out_channels))
|
||||
|
||||
self.resnets = nn.ModuleList(resnets)
|
||||
|
||||
@@ -527,24 +559,28 @@ class DownBlock1DNoSkip(nn.Module):
|
||||
|
||||
|
||||
class AttnUpBlock1D(nn.Module):
|
||||
def __init__(self, in_channels, out_channels, mid_channels=None):
|
||||
def __init__(self, in_channels: int, out_channels: int, mid_channels: int = None, num_layers: int = 3):
|
||||
super().__init__()
|
||||
if num_layers < 1:
|
||||
raise ValueError("AttnUpBlock1D requires added num_layers >= 1")
|
||||
|
||||
mid_channels = out_channels if mid_channels is None else mid_channels
|
||||
|
||||
resnets = [
|
||||
ResConvBlock(2 * in_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, out_channels),
|
||||
]
|
||||
attentions = [
|
||||
SelfAttention1d(mid_channels, mid_channels // 32),
|
||||
SelfAttention1d(mid_channels, mid_channels // 32),
|
||||
SelfAttention1d(out_channels, out_channels // 32),
|
||||
]
|
||||
resnets = []
|
||||
attentions = []
|
||||
|
||||
for i in range(num_layers):
|
||||
in_channels = 2 * in_channels if i == 0 else mid_channels
|
||||
if i < (num_layers - 1):
|
||||
resnets.append(ResConvBlock(in_channels, mid_channels, mid_channels))
|
||||
attentions.append(SelfAttention1d(mid_channels, mid_channels // 32))
|
||||
else:
|
||||
resnets.append(ResConvBlock(mid_channels, mid_channels, out_channels))
|
||||
attentions.append(SelfAttention1d(out_channels, out_channels // 32))
|
||||
|
||||
self.attentions = nn.ModuleList(attentions)
|
||||
self.resnets = nn.ModuleList(resnets)
|
||||
self.up = Upsample1d(kernel="cubic")
|
||||
self.up = KernelUpsample1D(kernel="cubic")
|
||||
|
||||
def forward(self, hidden_states, res_hidden_states_tuple, temb=None):
|
||||
res_hidden_states = res_hidden_states_tuple[-1]
|
||||
@@ -560,18 +596,24 @@ class AttnUpBlock1D(nn.Module):
|
||||
|
||||
|
||||
class UpBlock1D(nn.Module):
|
||||
def __init__(self, in_channels, out_channels, mid_channels=None):
|
||||
def __init__(self, in_channels: int, out_channels: int, mid_channels: int = None, num_layers: int = 3):
|
||||
super().__init__()
|
||||
if num_layers < 1:
|
||||
raise ValueError("UpBlock1D requires added num_layers >= 1")
|
||||
|
||||
mid_channels = in_channels if mid_channels is None else mid_channels
|
||||
|
||||
resnets = [
|
||||
ResConvBlock(2 * in_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, out_channels),
|
||||
]
|
||||
resnets = []
|
||||
|
||||
for i in range(num_layers):
|
||||
in_channels = 2 * in_channels if i == 0 else mid_channels
|
||||
if i < (num_layers - 1):
|
||||
resnets.append(ResConvBlock(in_channels, mid_channels, mid_channels))
|
||||
else:
|
||||
resnets.append(ResConvBlock(mid_channels, mid_channels, out_channels))
|
||||
|
||||
self.resnets = nn.ModuleList(resnets)
|
||||
self.up = Upsample1d(kernel="cubic")
|
||||
self.up = KernelUpsample1D(kernel="cubic")
|
||||
|
||||
def forward(self, hidden_states, res_hidden_states_tuple, temb=None):
|
||||
res_hidden_states = res_hidden_states_tuple[-1]
|
||||
@@ -586,15 +628,21 @@ class UpBlock1D(nn.Module):
|
||||
|
||||
|
||||
class UpBlock1DNoSkip(nn.Module):
|
||||
def __init__(self, in_channels, out_channels, mid_channels=None):
|
||||
def __init__(self, in_channels: int, out_channels: int, mid_channels: int = None, num_layers: int = 3):
|
||||
super().__init__()
|
||||
if num_layers < 1:
|
||||
raise ValueError("UpBlock1D requires added num_layers >= 1")
|
||||
|
||||
mid_channels = in_channels if mid_channels is None else mid_channels
|
||||
|
||||
resnets = [
|
||||
ResConvBlock(2 * in_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, mid_channels),
|
||||
ResConvBlock(mid_channels, mid_channels, out_channels, is_last=True),
|
||||
]
|
||||
resnets = []
|
||||
|
||||
for i in range(num_layers):
|
||||
in_channels = 2 * in_channels if i == 0 else mid_channels
|
||||
if i < (num_layers - 1):
|
||||
resnets.append(ResConvBlock(in_channels, mid_channels, mid_channels))
|
||||
else:
|
||||
resnets.append(ResConvBlock(mid_channels, mid_channels, out_channels, is_last=True))
|
||||
|
||||
self.resnets = nn.ModuleList(resnets)
|
||||
|
||||
|
||||
@@ -34,6 +34,7 @@ def get_down_block(
|
||||
downsample_padding=None,
|
||||
dual_cross_attention=False,
|
||||
use_linear_projection=False,
|
||||
only_cross_attention=False,
|
||||
):
|
||||
down_block_type = down_block_type[7:] if down_block_type.startswith("UNetRes") else down_block_type
|
||||
if down_block_type == "DownBlock2D":
|
||||
@@ -78,6 +79,7 @@ def get_down_block(
|
||||
attn_num_head_channels=attn_num_head_channels,
|
||||
dual_cross_attention=dual_cross_attention,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention,
|
||||
)
|
||||
elif down_block_type == "SkipDownBlock2D":
|
||||
return SkipDownBlock2D(
|
||||
@@ -143,6 +145,7 @@ def get_up_block(
|
||||
cross_attention_dim=None,
|
||||
dual_cross_attention=False,
|
||||
use_linear_projection=False,
|
||||
only_cross_attention=False,
|
||||
):
|
||||
up_block_type = up_block_type[7:] if up_block_type.startswith("UNetRes") else up_block_type
|
||||
if up_block_type == "UpBlock2D":
|
||||
@@ -174,6 +177,7 @@ def get_up_block(
|
||||
attn_num_head_channels=attn_num_head_channels,
|
||||
dual_cross_attention=dual_cross_attention,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention,
|
||||
)
|
||||
elif up_block_type == "AttnUpBlock2D":
|
||||
return AttnUpBlock2D(
|
||||
@@ -250,7 +254,6 @@ class UNetMidBlock2D(nn.Module):
|
||||
attn_num_head_channels=1,
|
||||
attention_type="default",
|
||||
output_scale_factor=1.0,
|
||||
**kwargs,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -332,7 +335,6 @@ class UNetMidBlock2DCrossAttn(nn.Module):
|
||||
cross_attention_dim=1280,
|
||||
dual_cross_attention=False,
|
||||
use_linear_projection=False,
|
||||
**kwargs,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -400,15 +402,17 @@ class UNetMidBlock2DCrossAttn(nn.Module):
|
||||
self.resnets = nn.ModuleList(resnets)
|
||||
|
||||
def set_attention_slice(self, slice_size):
|
||||
if slice_size is not None and self.attn_num_head_channels % slice_size != 0:
|
||||
head_dims = self.attn_num_head_channels
|
||||
head_dims = [head_dims] if isinstance(head_dims, int) else head_dims
|
||||
if slice_size is not None and any(dim % slice_size != 0 for dim in head_dims):
|
||||
raise ValueError(
|
||||
f"Make sure slice_size {slice_size} is a divisor of "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"Make sure slice_size {slice_size} is a common divisor of "
|
||||
f"the number of heads used in cross_attention: {head_dims}"
|
||||
)
|
||||
if slice_size is not None and slice_size > self.attn_num_head_channels:
|
||||
if slice_size is not None and slice_size > min(head_dims):
|
||||
raise ValueError(
|
||||
f"Chunk_size {slice_size} has to be smaller or equal to "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"slice_size {slice_size} has to be smaller or equal to "
|
||||
f"the lowest number of heads used in cross_attention: min({head_dims}) = {min(head_dims)}"
|
||||
)
|
||||
|
||||
for attn in self.attentions:
|
||||
@@ -530,6 +534,7 @@ class CrossAttnDownBlock2D(nn.Module):
|
||||
add_downsample=True,
|
||||
dual_cross_attention=False,
|
||||
use_linear_projection=False,
|
||||
only_cross_attention=False,
|
||||
):
|
||||
super().__init__()
|
||||
resnets = []
|
||||
@@ -564,6 +569,7 @@ class CrossAttnDownBlock2D(nn.Module):
|
||||
cross_attention_dim=cross_attention_dim,
|
||||
norm_num_groups=resnet_groups,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention,
|
||||
)
|
||||
)
|
||||
else:
|
||||
@@ -594,15 +600,17 @@ class CrossAttnDownBlock2D(nn.Module):
|
||||
self.gradient_checkpointing = False
|
||||
|
||||
def set_attention_slice(self, slice_size):
|
||||
if slice_size is not None and self.attn_num_head_channels % slice_size != 0:
|
||||
head_dims = self.attn_num_head_channels
|
||||
head_dims = [head_dims] if isinstance(head_dims, int) else head_dims
|
||||
if slice_size is not None and any(dim % slice_size != 0 for dim in head_dims):
|
||||
raise ValueError(
|
||||
f"Make sure slice_size {slice_size} is a divisor of "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"Make sure slice_size {slice_size} is a common divisor of "
|
||||
f"the number of heads used in cross_attention: {head_dims}"
|
||||
)
|
||||
if slice_size is not None and slice_size > self.attn_num_head_channels:
|
||||
if slice_size is not None and slice_size > min(head_dims):
|
||||
raise ValueError(
|
||||
f"Chunk_size {slice_size} has to be smaller or equal to "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"slice_size {slice_size} has to be smaller or equal to "
|
||||
f"the lowest number of heads used in cross_attention: min({head_dims}) = {min(head_dims)}"
|
||||
)
|
||||
|
||||
for attn in self.attentions:
|
||||
@@ -1129,6 +1137,7 @@ class CrossAttnUpBlock2D(nn.Module):
|
||||
add_upsample=True,
|
||||
dual_cross_attention=False,
|
||||
use_linear_projection=False,
|
||||
only_cross_attention=False,
|
||||
):
|
||||
super().__init__()
|
||||
resnets = []
|
||||
@@ -1165,6 +1174,7 @@ class CrossAttnUpBlock2D(nn.Module):
|
||||
cross_attention_dim=cross_attention_dim,
|
||||
norm_num_groups=resnet_groups,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention,
|
||||
)
|
||||
)
|
||||
else:
|
||||
@@ -1189,15 +1199,17 @@ class CrossAttnUpBlock2D(nn.Module):
|
||||
self.gradient_checkpointing = False
|
||||
|
||||
def set_attention_slice(self, slice_size):
|
||||
if slice_size is not None and self.attn_num_head_channels % slice_size != 0:
|
||||
head_dims = self.attn_num_head_channels
|
||||
head_dims = [head_dims] if isinstance(head_dims, int) else head_dims
|
||||
if slice_size is not None and any(dim % slice_size != 0 for dim in head_dims):
|
||||
raise ValueError(
|
||||
f"Make sure slice_size {slice_size} is a divisor of "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"Make sure slice_size {slice_size} is a common divisor of "
|
||||
f"the number of heads used in cross_attention: {head_dims}"
|
||||
)
|
||||
if slice_size is not None and slice_size > self.attn_num_head_channels:
|
||||
if slice_size is not None and slice_size > min(head_dims):
|
||||
raise ValueError(
|
||||
f"Chunk_size {slice_size} has to be smaller or equal to "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"slice_size {slice_size} has to be smaller or equal to "
|
||||
f"the lowest number of heads used in cross_attention: min({head_dims}) = {min(head_dims)}"
|
||||
)
|
||||
|
||||
for attn in self.attentions:
|
||||
|
||||
@@ -46,6 +46,8 @@ class FlaxCrossAttnDownBlock2D(nn.Module):
|
||||
num_layers: int = 1
|
||||
attn_num_head_channels: int = 1
|
||||
add_downsample: bool = True
|
||||
use_linear_projection: bool = False
|
||||
only_cross_attention: bool = False
|
||||
dtype: jnp.dtype = jnp.float32
|
||||
|
||||
def setup(self):
|
||||
@@ -68,6 +70,8 @@ class FlaxCrossAttnDownBlock2D(nn.Module):
|
||||
n_heads=self.attn_num_head_channels,
|
||||
d_head=self.out_channels // self.attn_num_head_channels,
|
||||
depth=1,
|
||||
use_linear_projection=self.use_linear_projection,
|
||||
only_cross_attention=self.only_cross_attention,
|
||||
dtype=self.dtype,
|
||||
)
|
||||
attentions.append(attn_block)
|
||||
@@ -178,6 +182,8 @@ class FlaxCrossAttnUpBlock2D(nn.Module):
|
||||
num_layers: int = 1
|
||||
attn_num_head_channels: int = 1
|
||||
add_upsample: bool = True
|
||||
use_linear_projection: bool = False
|
||||
only_cross_attention: bool = False
|
||||
dtype: jnp.dtype = jnp.float32
|
||||
|
||||
def setup(self):
|
||||
@@ -201,6 +207,8 @@ class FlaxCrossAttnUpBlock2D(nn.Module):
|
||||
n_heads=self.attn_num_head_channels,
|
||||
d_head=self.out_channels // self.attn_num_head_channels,
|
||||
depth=1,
|
||||
use_linear_projection=self.use_linear_projection,
|
||||
only_cross_attention=self.only_cross_attention,
|
||||
dtype=self.dtype,
|
||||
)
|
||||
attentions.append(attn_block)
|
||||
@@ -310,6 +318,7 @@ class FlaxUNetMidBlock2DCrossAttn(nn.Module):
|
||||
dropout: float = 0.0
|
||||
num_layers: int = 1
|
||||
attn_num_head_channels: int = 1
|
||||
use_linear_projection: bool = False
|
||||
dtype: jnp.dtype = jnp.float32
|
||||
|
||||
def setup(self):
|
||||
@@ -331,6 +340,7 @@ class FlaxUNetMidBlock2DCrossAttn(nn.Module):
|
||||
n_heads=self.attn_num_head_channels,
|
||||
d_head=self.in_channels // self.attn_num_head_channels,
|
||||
depth=1,
|
||||
use_linear_projection=self.use_linear_projection,
|
||||
dtype=self.dtype,
|
||||
)
|
||||
attentions.append(attn_block)
|
||||
|
||||
@@ -98,6 +98,7 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin):
|
||||
"DownBlock2D",
|
||||
),
|
||||
up_block_types: Tuple[str] = ("UpBlock2D", "CrossAttnUpBlock2D", "CrossAttnUpBlock2D", "CrossAttnUpBlock2D"),
|
||||
only_cross_attention: Union[bool, Tuple[bool]] = False,
|
||||
block_out_channels: Tuple[int] = (320, 640, 1280, 1280),
|
||||
layers_per_block: int = 2,
|
||||
downsample_padding: int = 1,
|
||||
@@ -109,6 +110,7 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin):
|
||||
attention_head_dim: Union[int, Tuple[int]] = 8,
|
||||
dual_cross_attention: bool = False,
|
||||
use_linear_projection: bool = False,
|
||||
num_class_embeds: Optional[int] = None,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -124,10 +126,17 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin):
|
||||
|
||||
self.time_embedding = TimestepEmbedding(timestep_input_dim, time_embed_dim)
|
||||
|
||||
# class embedding
|
||||
if num_class_embeds is not None:
|
||||
self.class_embedding = nn.Embedding(num_class_embeds, time_embed_dim)
|
||||
|
||||
self.down_blocks = nn.ModuleList([])
|
||||
self.mid_block = None
|
||||
self.up_blocks = nn.ModuleList([])
|
||||
|
||||
if isinstance(only_cross_attention, bool):
|
||||
only_cross_attention = [only_cross_attention] * len(down_block_types)
|
||||
|
||||
if isinstance(attention_head_dim, int):
|
||||
attention_head_dim = (attention_head_dim,) * len(down_block_types)
|
||||
|
||||
@@ -153,6 +162,7 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin):
|
||||
downsample_padding=downsample_padding,
|
||||
dual_cross_attention=dual_cross_attention,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention[i],
|
||||
)
|
||||
self.down_blocks.append(down_block)
|
||||
|
||||
@@ -177,6 +187,7 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin):
|
||||
# up
|
||||
reversed_block_out_channels = list(reversed(block_out_channels))
|
||||
reversed_attention_head_dim = list(reversed(attention_head_dim))
|
||||
only_cross_attention = list(reversed(only_cross_attention))
|
||||
output_channel = reversed_block_out_channels[0]
|
||||
for i, up_block_type in enumerate(up_block_types):
|
||||
is_final_block = i == len(block_out_channels) - 1
|
||||
@@ -207,6 +218,7 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin):
|
||||
attn_num_head_channels=reversed_attention_head_dim[i],
|
||||
dual_cross_attention=dual_cross_attention,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention[i],
|
||||
)
|
||||
self.up_blocks.append(up_block)
|
||||
prev_output_channel = output_channel
|
||||
@@ -217,15 +229,17 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin):
|
||||
self.conv_out = nn.Conv2d(block_out_channels[0], out_channels, kernel_size=3, padding=1)
|
||||
|
||||
def set_attention_slice(self, slice_size):
|
||||
if slice_size is not None and self.config.attention_head_dim % slice_size != 0:
|
||||
head_dims = self.config.attention_head_dim
|
||||
head_dims = [head_dims] if isinstance(head_dims, int) else head_dims
|
||||
if slice_size is not None and any(dim % slice_size != 0 for dim in head_dims):
|
||||
raise ValueError(
|
||||
f"Make sure slice_size {slice_size} is a divisor of "
|
||||
f"the number of heads used in cross_attention {self.config.attention_head_dim}"
|
||||
f"Make sure slice_size {slice_size} is a common divisor of "
|
||||
f"the number of heads used in cross_attention: {head_dims}"
|
||||
)
|
||||
if slice_size is not None and slice_size > self.config.attention_head_dim:
|
||||
if slice_size is not None and slice_size > min(head_dims):
|
||||
raise ValueError(
|
||||
f"Chunk_size {slice_size} has to be smaller or equal to "
|
||||
f"the number of heads used in cross_attention {self.config.attention_head_dim}"
|
||||
f"slice_size {slice_size} has to be smaller or equal to "
|
||||
f"the lowest number of heads used in cross_attention: min({head_dims}) = {min(head_dims)}"
|
||||
)
|
||||
|
||||
for block in self.down_blocks:
|
||||
@@ -258,6 +272,7 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin):
|
||||
sample: torch.FloatTensor,
|
||||
timestep: Union[torch.Tensor, float, int],
|
||||
encoder_hidden_states: torch.Tensor,
|
||||
class_labels: Optional[torch.Tensor] = None,
|
||||
return_dict: bool = True,
|
||||
) -> Union[UNet2DConditionOutput, Tuple]:
|
||||
r"""
|
||||
@@ -310,6 +325,12 @@ class UNet2DConditionModel(ModelMixin, ConfigMixin):
|
||||
t_emb = t_emb.to(dtype=self.dtype)
|
||||
emb = self.time_embedding(t_emb)
|
||||
|
||||
if self.config.num_class_embeds is not None:
|
||||
if class_labels is None:
|
||||
raise ValueError("class_labels should be provided when num_class_embeds > 0")
|
||||
class_emb = self.class_embedding(class_labels).to(dtype=self.dtype)
|
||||
emb = emb + class_emb
|
||||
|
||||
# 2. pre-process
|
||||
sample = self.conv_in(sample)
|
||||
|
||||
|
||||
@@ -79,7 +79,7 @@ class FlaxUNet2DConditionModel(nn.Module, FlaxModelMixin, ConfigMixin):
|
||||
The tuple of output channels for each block.
|
||||
layers_per_block (`int`, *optional*, defaults to 2):
|
||||
The number of layers per block.
|
||||
attention_head_dim (`int`, *optional*, defaults to 8):
|
||||
attention_head_dim (`int` or `Tuple[int]`, *optional*, defaults to 8):
|
||||
The dimension of the attention heads.
|
||||
cross_attention_dim (`int`, *optional*, defaults to 768):
|
||||
The dimension of the cross attention features.
|
||||
@@ -97,11 +97,13 @@ class FlaxUNet2DConditionModel(nn.Module, FlaxModelMixin, ConfigMixin):
|
||||
"DownBlock2D",
|
||||
)
|
||||
up_block_types: Tuple[str] = ("UpBlock2D", "CrossAttnUpBlock2D", "CrossAttnUpBlock2D", "CrossAttnUpBlock2D")
|
||||
only_cross_attention: Union[bool, Tuple[bool]] = False
|
||||
block_out_channels: Tuple[int] = (320, 640, 1280, 1280)
|
||||
layers_per_block: int = 2
|
||||
attention_head_dim: int = 8
|
||||
attention_head_dim: Union[int, Tuple[int]] = 8
|
||||
cross_attention_dim: int = 1280
|
||||
dropout: float = 0.0
|
||||
use_linear_projection: bool = False
|
||||
dtype: jnp.dtype = jnp.float32
|
||||
freq_shift: int = 0
|
||||
|
||||
@@ -134,6 +136,14 @@ class FlaxUNet2DConditionModel(nn.Module, FlaxModelMixin, ConfigMixin):
|
||||
self.time_proj = FlaxTimesteps(block_out_channels[0], freq_shift=self.config.freq_shift)
|
||||
self.time_embedding = FlaxTimestepEmbedding(time_embed_dim, dtype=self.dtype)
|
||||
|
||||
only_cross_attention = self.only_cross_attention
|
||||
if isinstance(only_cross_attention, bool):
|
||||
only_cross_attention = (only_cross_attention,) * len(self.down_block_types)
|
||||
|
||||
attention_head_dim = self.attention_head_dim
|
||||
if isinstance(attention_head_dim, int):
|
||||
attention_head_dim = (attention_head_dim,) * len(self.down_block_types)
|
||||
|
||||
# down
|
||||
down_blocks = []
|
||||
output_channel = block_out_channels[0]
|
||||
@@ -148,8 +158,10 @@ class FlaxUNet2DConditionModel(nn.Module, FlaxModelMixin, ConfigMixin):
|
||||
out_channels=output_channel,
|
||||
dropout=self.dropout,
|
||||
num_layers=self.layers_per_block,
|
||||
attn_num_head_channels=self.attention_head_dim,
|
||||
attn_num_head_channels=attention_head_dim[i],
|
||||
add_downsample=not is_final_block,
|
||||
use_linear_projection=self.use_linear_projection,
|
||||
only_cross_attention=only_cross_attention[i],
|
||||
dtype=self.dtype,
|
||||
)
|
||||
else:
|
||||
@@ -169,13 +181,16 @@ class FlaxUNet2DConditionModel(nn.Module, FlaxModelMixin, ConfigMixin):
|
||||
self.mid_block = FlaxUNetMidBlock2DCrossAttn(
|
||||
in_channels=block_out_channels[-1],
|
||||
dropout=self.dropout,
|
||||
attn_num_head_channels=self.attention_head_dim,
|
||||
attn_num_head_channels=attention_head_dim[-1],
|
||||
use_linear_projection=self.use_linear_projection,
|
||||
dtype=self.dtype,
|
||||
)
|
||||
|
||||
# up
|
||||
up_blocks = []
|
||||
reversed_block_out_channels = list(reversed(block_out_channels))
|
||||
reversed_attention_head_dim = list(reversed(attention_head_dim))
|
||||
only_cross_attention = list(reversed(only_cross_attention))
|
||||
output_channel = reversed_block_out_channels[0]
|
||||
for i, up_block_type in enumerate(self.up_block_types):
|
||||
prev_output_channel = output_channel
|
||||
@@ -190,9 +205,11 @@ class FlaxUNet2DConditionModel(nn.Module, FlaxModelMixin, ConfigMixin):
|
||||
out_channels=output_channel,
|
||||
prev_output_channel=prev_output_channel,
|
||||
num_layers=self.layers_per_block + 1,
|
||||
attn_num_head_channels=self.attention_head_dim,
|
||||
attn_num_head_channels=reversed_attention_head_dim[i],
|
||||
add_upsample=not is_final_block,
|
||||
dropout=self.dropout,
|
||||
use_linear_projection=self.use_linear_projection,
|
||||
only_cross_attention=only_cross_attention[i],
|
||||
dtype=self.dtype,
|
||||
)
|
||||
else:
|
||||
|
||||
@@ -565,6 +565,7 @@ class AutoencoderKL(ModelMixin, ConfigMixin):
|
||||
|
||||
self.quant_conv = torch.nn.Conv2d(2 * latent_channels, 2 * latent_channels, 1)
|
||||
self.post_quant_conv = torch.nn.Conv2d(latent_channels, latent_channels, 1)
|
||||
self.use_slicing = False
|
||||
|
||||
def encode(self, x: torch.FloatTensor, return_dict: bool = True) -> AutoencoderKLOutput:
|
||||
h = self.encoder(x)
|
||||
@@ -576,7 +577,7 @@ class AutoencoderKL(ModelMixin, ConfigMixin):
|
||||
|
||||
return AutoencoderKLOutput(latent_dist=posterior)
|
||||
|
||||
def decode(self, z: torch.FloatTensor, return_dict: bool = True) -> Union[DecoderOutput, torch.FloatTensor]:
|
||||
def _decode(self, z: torch.FloatTensor, return_dict: bool = True) -> Union[DecoderOutput, torch.FloatTensor]:
|
||||
z = self.post_quant_conv(z)
|
||||
dec = self.decoder(z)
|
||||
|
||||
@@ -585,6 +586,34 @@ class AutoencoderKL(ModelMixin, ConfigMixin):
|
||||
|
||||
return DecoderOutput(sample=dec)
|
||||
|
||||
def enable_slicing(self):
|
||||
r"""
|
||||
Enable sliced VAE decoding.
|
||||
|
||||
When this option is enabled, the VAE will split the input tensor in slices to compute decoding in several
|
||||
steps. This is useful to save some memory and allow larger batch sizes.
|
||||
"""
|
||||
self.use_slicing = True
|
||||
|
||||
def disable_slicing(self):
|
||||
r"""
|
||||
Disable sliced VAE decoding. If `enable_slicing` was previously invoked, this method will go back to computing
|
||||
decoding in one step.
|
||||
"""
|
||||
self.use_slicing = False
|
||||
|
||||
def decode(self, z: torch.FloatTensor, return_dict: bool = True) -> Union[DecoderOutput, torch.FloatTensor]:
|
||||
if self.use_slicing and z.shape[0] > 1:
|
||||
decoded_slices = [self._decode(z_slice).sample for z_slice in z.split(1)]
|
||||
decoded = torch.cat(decoded_slices)
|
||||
else:
|
||||
decoded = self._decode(z).sample
|
||||
|
||||
if not return_dict:
|
||||
return (decoded,)
|
||||
|
||||
return DecoderOutput(sample=decoded)
|
||||
|
||||
def forward(
|
||||
self,
|
||||
sample: torch.FloatTensor,
|
||||
|
||||
@@ -317,8 +317,8 @@ class FlaxDiffusionPipeline(ConfigMixin):
|
||||
allow_patterns = [os.path.join(k, "*") for k in folder_names]
|
||||
allow_patterns += [FLAX_WEIGHTS_NAME, SCHEDULER_CONFIG_NAME, CONFIG_NAME, cls.config_name]
|
||||
|
||||
# make sure we don't download PyTorch weights
|
||||
ignore_patterns = "*.bin"
|
||||
# make sure we don't download PyTorch weights, unless when using from_pt
|
||||
ignore_patterns = "*.bin" if not from_pt else []
|
||||
|
||||
if cls != FlaxDiffusionPipeline:
|
||||
requested_pipeline_class = cls.__name__
|
||||
|
||||
@@ -26,7 +26,7 @@ import torch
|
||||
|
||||
import diffusers
|
||||
import PIL
|
||||
from huggingface_hub import snapshot_download
|
||||
from huggingface_hub import model_info, snapshot_download
|
||||
from packaging import version
|
||||
from PIL import Image
|
||||
from tqdm.auto import tqdm
|
||||
@@ -44,6 +44,7 @@ from .utils import (
|
||||
BaseOutput,
|
||||
deprecate,
|
||||
is_accelerate_available,
|
||||
is_safetensors_available,
|
||||
is_torch_version,
|
||||
is_transformers_available,
|
||||
logging,
|
||||
@@ -117,6 +118,23 @@ class AudioPipelineOutput(BaseOutput):
|
||||
audios: np.ndarray
|
||||
|
||||
|
||||
def is_safetensors_compatible(info) -> bool:
|
||||
filenames = set(sibling.rfilename for sibling in info.siblings)
|
||||
pt_filenames = set(filename for filename in filenames if filename.endswith(".bin"))
|
||||
is_safetensors_compatible = any(file.endswith(".safetensors") for file in filenames)
|
||||
for pt_filename in pt_filenames:
|
||||
prefix, raw = os.path.split(pt_filename)
|
||||
if raw == "pytorch_model.bin":
|
||||
# transformers specific
|
||||
sf_filename = os.path.join(prefix, "model.safetensors")
|
||||
else:
|
||||
sf_filename = pt_filename[: -len(".bin")] + ".safetensors"
|
||||
if is_safetensors_compatible and sf_filename not in filenames:
|
||||
logger.warning(f"{sf_filename} not found")
|
||||
is_safetensors_compatible = False
|
||||
return is_safetensors_compatible
|
||||
|
||||
|
||||
class DiffusionPipeline(ConfigMixin):
|
||||
r"""
|
||||
Base class for all models.
|
||||
@@ -129,10 +147,13 @@ class DiffusionPipeline(ConfigMixin):
|
||||
|
||||
Class attributes:
|
||||
|
||||
- **config_name** ([`str`]) -- name of the config file that will store the class and module names of all
|
||||
- **config_name** (`str`) -- name of the config file that will store the class and module names of all
|
||||
components of the diffusion pipeline.
|
||||
- **_optional_components** (List[`str`]) -- list of all components that are optional so they don't have to be
|
||||
passed for the pipeline to function (should be overridden by subclasses).
|
||||
"""
|
||||
config_name = "model_index.json"
|
||||
_optional_components = []
|
||||
|
||||
def register_modules(self, **kwargs):
|
||||
# import it here to avoid circular import
|
||||
@@ -184,12 +205,19 @@ class DiffusionPipeline(ConfigMixin):
|
||||
model_index_dict.pop("_diffusers_version")
|
||||
model_index_dict.pop("_module", None)
|
||||
|
||||
expected_modules, optional_kwargs = self._get_signature_keys(self)
|
||||
|
||||
def is_saveable_module(name, value):
|
||||
if name not in expected_modules:
|
||||
return False
|
||||
if name in self._optional_components and value[0] is None:
|
||||
return False
|
||||
return True
|
||||
|
||||
model_index_dict = {k: v for k, v in model_index_dict.items() if is_saveable_module(k, v)}
|
||||
|
||||
for pipeline_component_name in model_index_dict.keys():
|
||||
sub_model = getattr(self, pipeline_component_name)
|
||||
if sub_model is None:
|
||||
# edge case for saving a pipeline with safety_checker=None
|
||||
continue
|
||||
|
||||
model_cls = sub_model.__class__
|
||||
|
||||
save_method_name = None
|
||||
@@ -449,7 +477,7 @@ class DiffusionPipeline(ConfigMixin):
|
||||
allow_patterns += [WEIGHTS_NAME, SCHEDULER_CONFIG_NAME, CONFIG_NAME, ONNX_WEIGHTS_NAME, cls.config_name]
|
||||
|
||||
# make sure we don't download flax weights
|
||||
ignore_patterns = "*.msgpack"
|
||||
ignore_patterns = ["*.msgpack"]
|
||||
|
||||
if custom_pipeline is not None:
|
||||
allow_patterns += [CUSTOM_PIPELINE_FILE_NAME]
|
||||
@@ -463,6 +491,15 @@ class DiffusionPipeline(ConfigMixin):
|
||||
user_agent["custom_pipeline"] = custom_pipeline
|
||||
user_agent = http_user_agent(user_agent)
|
||||
|
||||
if is_safetensors_available():
|
||||
info = model_info(
|
||||
pretrained_model_name_or_path,
|
||||
use_auth_token=use_auth_token,
|
||||
revision=revision,
|
||||
)
|
||||
if is_safetensors_compatible(info):
|
||||
ignore_patterns.append("*.bin")
|
||||
|
||||
# download all allow_patterns
|
||||
cached_folder = snapshot_download(
|
||||
pretrained_model_name_or_path,
|
||||
@@ -523,38 +560,47 @@ class DiffusionPipeline(ConfigMixin):
|
||||
# some modules can be passed directly to the init
|
||||
# in this case they are already instantiated in `kwargs`
|
||||
# extract them here
|
||||
expected_modules = set(inspect.signature(pipeline_class.__init__).parameters.keys()) - set(["self"])
|
||||
expected_modules, optional_kwargs = cls._get_signature_keys(pipeline_class)
|
||||
passed_class_obj = {k: kwargs.pop(k) for k in expected_modules if k in kwargs}
|
||||
passed_pipe_kwargs = {k: kwargs.pop(k) for k in optional_kwargs if k in kwargs}
|
||||
|
||||
init_dict, unused_kwargs, _ = pipeline_class.extract_init_dict(config_dict, **kwargs)
|
||||
|
||||
if len(unused_kwargs) > 0:
|
||||
logger.warning(f"Keyword arguments {unused_kwargs} not recognized.")
|
||||
# define init kwargs
|
||||
init_kwargs = {k: init_dict.pop(k) for k in optional_kwargs if k in init_dict}
|
||||
init_kwargs = {**init_kwargs, **passed_pipe_kwargs}
|
||||
|
||||
init_kwargs = {}
|
||||
# remove `null` components
|
||||
def load_module(name, value):
|
||||
if value[0] is None:
|
||||
return False
|
||||
if name in passed_class_obj and passed_class_obj[name] is None:
|
||||
return False
|
||||
return True
|
||||
|
||||
init_dict = {k: v for k, v in init_dict.items() if load_module(k, v)}
|
||||
|
||||
if len(unused_kwargs) > 0:
|
||||
logger.warning(
|
||||
f"Keyword arguments {unused_kwargs} are not expected by {pipeline_class.__name__} and will be ignored."
|
||||
)
|
||||
|
||||
# import it here to avoid circular import
|
||||
from diffusers import pipelines
|
||||
|
||||
# 3. Load each module in the pipeline
|
||||
for name, (library_name, class_name) in init_dict.items():
|
||||
if class_name is None:
|
||||
# edge case for when the pipeline was saved with safety_checker=None
|
||||
init_kwargs[name] = None
|
||||
continue
|
||||
|
||||
# 3.1 - now that JAX/Flax is an official framework of the library, we might load from Flax names
|
||||
if class_name.startswith("Flax"):
|
||||
class_name = class_name[4:]
|
||||
|
||||
is_pipeline_module = hasattr(pipelines, library_name)
|
||||
loaded_sub_model = None
|
||||
sub_model_should_be_defined = True
|
||||
|
||||
# if the model is in a pipeline module, then we load it from the pipeline
|
||||
if name in passed_class_obj:
|
||||
# 1. check that passed_class_obj has correct parent class
|
||||
if not is_pipeline_module and passed_class_obj[name] is not None:
|
||||
if not is_pipeline_module:
|
||||
library = importlib.import_module(library_name)
|
||||
class_obj = getattr(library, class_name)
|
||||
importable_classes = LOADABLE_CLASSES[library_name]
|
||||
@@ -570,12 +616,6 @@ class DiffusionPipeline(ConfigMixin):
|
||||
f"{passed_class_obj[name]} is of type: {type(passed_class_obj[name])}, but should be"
|
||||
f" {expected_class_obj}"
|
||||
)
|
||||
elif passed_class_obj[name] is None:
|
||||
logger.warning(
|
||||
f"You have passed `None` for {name} to disable its functionality in {pipeline_class}. Note"
|
||||
f" that this might lead to problems when using {pipeline_class} and is not recommended."
|
||||
)
|
||||
sub_model_should_be_defined = False
|
||||
else:
|
||||
logger.warning(
|
||||
f"You have passed a non-standard module {passed_class_obj[name]}. We cannot verify whether it"
|
||||
@@ -597,7 +637,7 @@ class DiffusionPipeline(ConfigMixin):
|
||||
importable_classes = LOADABLE_CLASSES[library_name]
|
||||
class_candidates = {c: getattr(library, c, None) for c in importable_classes.keys()}
|
||||
|
||||
if loaded_sub_model is None and sub_model_should_be_defined:
|
||||
if loaded_sub_model is None:
|
||||
load_method_name = None
|
||||
for class_name, class_candidate in class_candidates.items():
|
||||
if class_candidate is not None and issubclass(class_obj, class_candidate):
|
||||
@@ -651,11 +691,13 @@ class DiffusionPipeline(ConfigMixin):
|
||||
|
||||
# 4. Potentially add passed objects if expected
|
||||
missing_modules = set(expected_modules) - set(init_kwargs.keys())
|
||||
if len(missing_modules) > 0 and missing_modules <= set(passed_class_obj.keys()):
|
||||
passed_modules = list(passed_class_obj.keys())
|
||||
optional_modules = pipeline_class._optional_components
|
||||
if len(missing_modules) > 0 and missing_modules <= set(passed_modules + optional_modules):
|
||||
for module in missing_modules:
|
||||
init_kwargs[module] = passed_class_obj[module]
|
||||
init_kwargs[module] = passed_class_obj.get(module, None)
|
||||
elif len(missing_modules) > 0:
|
||||
passed_modules = set(list(init_kwargs.keys()) + list(passed_class_obj.keys()))
|
||||
passed_modules = set(list(init_kwargs.keys()) + list(passed_class_obj.keys())) - optional_kwargs
|
||||
raise ValueError(
|
||||
f"Pipeline {pipeline_class} expected {expected_modules}, but only {passed_modules} were passed."
|
||||
)
|
||||
@@ -664,6 +706,14 @@ class DiffusionPipeline(ConfigMixin):
|
||||
model = pipeline_class(**init_kwargs)
|
||||
return model
|
||||
|
||||
@staticmethod
|
||||
def _get_signature_keys(obj):
|
||||
parameters = inspect.signature(obj.__init__).parameters
|
||||
required_parameters = {k: v for k, v in parameters.items() if v.default == inspect._empty}
|
||||
optional_parameters = set({k for k, v in parameters.items() if v.default != inspect._empty})
|
||||
expected_modules = set(required_parameters.keys()) - set(["self"])
|
||||
return expected_modules, optional_parameters
|
||||
|
||||
@property
|
||||
def components(self) -> Dict[str, Any]:
|
||||
r"""
|
||||
@@ -688,8 +738,10 @@ class DiffusionPipeline(ConfigMixin):
|
||||
Returns:
|
||||
A dictionaly containing all the modules needed to initialize the pipeline.
|
||||
"""
|
||||
components = {k: getattr(self, k) for k in self.config.keys() if not k.startswith("_")}
|
||||
expected_modules = set(inspect.signature(self.__init__).parameters.keys()) - set(["self"])
|
||||
expected_modules, optional_parameters = self._get_signature_keys(self)
|
||||
components = {
|
||||
k: getattr(self, k) for k in self.config.keys() if not k.startswith("_") and k not in optional_parameters
|
||||
}
|
||||
|
||||
if set(components.keys()) != expected_modules:
|
||||
raise ValueError(
|
||||
@@ -715,7 +767,7 @@ class DiffusionPipeline(ConfigMixin):
|
||||
|
||||
return pil_images
|
||||
|
||||
def progress_bar(self, iterable):
|
||||
def progress_bar(self, iterable=None, total=None):
|
||||
if not hasattr(self, "_progress_bar_config"):
|
||||
self._progress_bar_config = {}
|
||||
elif not isinstance(self._progress_bar_config, dict):
|
||||
@@ -723,7 +775,12 @@ class DiffusionPipeline(ConfigMixin):
|
||||
f"`self._progress_bar_config` should be of type `dict`, but is {type(self._progress_bar_config)}."
|
||||
)
|
||||
|
||||
return tqdm(iterable, **self._progress_bar_config)
|
||||
if iterable is not None:
|
||||
return tqdm(iterable, **self._progress_bar_config)
|
||||
elif total is not None:
|
||||
return tqdm(total=total, **self._progress_bar_config)
|
||||
else:
|
||||
raise ValueError("Either `total` or `iterable` has to be defined.")
|
||||
|
||||
def set_progress_bar_config(self, **kwargs):
|
||||
self._progress_bar_config = kwargs
|
||||
|
||||
@@ -126,7 +126,7 @@ init_image = init_image.resize((768, 512))
|
||||
|
||||
prompt = "A fantasy landscape, trending on artstation"
|
||||
|
||||
images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
|
||||
|
||||
images[0].save("fantasy_landscape.png")
|
||||
```
|
||||
|
||||
@@ -24,6 +24,7 @@ if is_torch_available() and is_transformers_available():
|
||||
StableDiffusionInpaintPipeline,
|
||||
StableDiffusionInpaintPipelineLegacy,
|
||||
StableDiffusionPipeline,
|
||||
StableDiffusionUpscalePipeline,
|
||||
)
|
||||
from .stable_diffusion_safe import StableDiffusionPipelineSafe
|
||||
from .versatile_diffusion import (
|
||||
|
||||
@@ -18,6 +18,7 @@ from typing import Callable, List, Optional, Union
|
||||
import torch
|
||||
|
||||
from diffusers.utils import is_accelerate_available
|
||||
from packaging import version
|
||||
from transformers import CLIPFeatureExtractor, XLMRobertaTokenizer
|
||||
|
||||
from ...configuration_utils import FrozenDict
|
||||
@@ -67,6 +68,7 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
feature_extractor ([`CLIPFeatureExtractor`]):
|
||||
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
|
||||
"""
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
@@ -84,6 +86,7 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
],
|
||||
safety_checker: StableDiffusionSafetyChecker,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -114,7 +117,7 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
new_config["clip_sample"] = False
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Alt Diffusion license and do not expose unfiltered"
|
||||
@@ -124,6 +127,33 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
|
||||
version.parse(unet.config._diffusers_version).base_version
|
||||
) < version.parse("0.9.0.dev0")
|
||||
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
|
||||
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
|
||||
deprecation_message = (
|
||||
"The configuration file of the unet has set the default `sample_size` to smaller than"
|
||||
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
|
||||
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
|
||||
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
|
||||
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
|
||||
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
|
||||
" in the config might lead to incorrect results in future versions. If you have downloaded this"
|
||||
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
|
||||
" the `unet/config.json` file"
|
||||
)
|
||||
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
|
||||
new_config = dict(unet.config)
|
||||
new_config["sample_size"] = 64
|
||||
unet._internal_dict = FrozenDict(new_config)
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
text_encoder=text_encoder,
|
||||
@@ -133,6 +163,8 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
def enable_xformers_memory_efficient_attention(self):
|
||||
r"""
|
||||
@@ -166,9 +198,14 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
if isinstance(self.unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.unet.config.attention_head_dim)
|
||||
|
||||
self.unet.set_attention_slice(slice_size)
|
||||
|
||||
def disable_attention_slicing(self):
|
||||
@@ -179,6 +216,22 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
# set slice_size = `None` to disable `attention slicing`
|
||||
self.enable_attention_slicing(None)
|
||||
|
||||
def enable_vae_slicing(self):
|
||||
r"""
|
||||
Enable sliced VAE decoding.
|
||||
|
||||
When this option is enabled, the VAE will split the input tensor in slices to compute decoding in several
|
||||
steps. This is useful to save some memory and allow larger batch sizes.
|
||||
"""
|
||||
self.vae.enable_slicing()
|
||||
|
||||
def disable_vae_slicing(self):
|
||||
r"""
|
||||
Disable sliced VAE decoding. If `enable_vae_slicing` was previously invoked, this method will go back to
|
||||
computing decoding in one step.
|
||||
"""
|
||||
self.vae.disable_slicing()
|
||||
|
||||
def enable_sequential_cpu_offload(self, gpu_id=0):
|
||||
r"""
|
||||
Offloads all models to CPU using accelerate, significantly reducing memory usage. When called, unet,
|
||||
@@ -192,10 +245,15 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
|
||||
device = torch.device(f"cuda:{gpu_id}")
|
||||
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae, self.safety_checker]:
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae]:
|
||||
if cpu_offloaded_model is not None:
|
||||
cpu_offload(cpu_offloaded_model, device)
|
||||
|
||||
if self.safety_checker is not None:
|
||||
# TODO(Patrick) - there is currently a bug with cpu offload of nn.Parameter in accelerate
|
||||
# fix by only offloading self.safety_checker for now
|
||||
cpu_offload(self.safety_checker.vision_model, device)
|
||||
|
||||
@property
|
||||
def _execution_device(self):
|
||||
r"""
|
||||
@@ -370,7 +428,7 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
)
|
||||
|
||||
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
|
||||
shape = (batch_size, num_channels_latents, height // 8, width // 8)
|
||||
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor)
|
||||
if latents is None:
|
||||
if device.type == "mps":
|
||||
# randn does not work reproducibly on mps
|
||||
@@ -390,8 +448,8 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -403,7 +461,6 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
return_dict: bool = True,
|
||||
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
|
||||
callback_steps: Optional[int] = 1,
|
||||
**kwargs,
|
||||
):
|
||||
r"""
|
||||
Function invoked when calling the pipeline for generation.
|
||||
@@ -411,9 +468,9 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -459,6 +516,9 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
# 1. Check inputs. Raise error if not correct
|
||||
self.check_inputs(prompt, height, width, callback_steps)
|
||||
@@ -497,25 +557,29 @@ class AltDiffusionPipeline(DiffusionPipeline):
|
||||
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
|
||||
|
||||
# 7. Denoising loop
|
||||
for i, t in enumerate(self.progress_bar(timesteps)):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
for i, t in enumerate(timesteps):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
|
||||
# call the callback, if provided
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
# call the callback, if provided
|
||||
if (i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0:
|
||||
progress_bar.update()
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
|
||||
# 8. Post-processing
|
||||
image = self.decode_latents(latents)
|
||||
|
||||
@@ -20,6 +20,7 @@ import torch
|
||||
|
||||
import PIL
|
||||
from diffusers.utils import is_accelerate_available
|
||||
from packaging import version
|
||||
from transformers import CLIPFeatureExtractor, XLMRobertaTokenizer
|
||||
|
||||
from ...configuration_utils import FrozenDict
|
||||
@@ -80,6 +81,7 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
feature_extractor ([`CLIPFeatureExtractor`]):
|
||||
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
|
||||
"""
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
@@ -97,6 +99,7 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
],
|
||||
safety_checker: StableDiffusionSafetyChecker,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -127,7 +130,7 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
new_config["clip_sample"] = False
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Alt Diffusion license and do not expose unfiltered"
|
||||
@@ -137,6 +140,33 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
|
||||
version.parse(unet.config._diffusers_version).base_version
|
||||
) < version.parse("0.9.0.dev0")
|
||||
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
|
||||
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
|
||||
deprecation_message = (
|
||||
"The configuration file of the unet has set the default `sample_size` to smaller than"
|
||||
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
|
||||
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
|
||||
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
|
||||
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
|
||||
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
|
||||
" in the config might lead to incorrect results in future versions. If you have downloaded this"
|
||||
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
|
||||
" the `unet/config.json` file"
|
||||
)
|
||||
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
|
||||
new_config = dict(unet.config)
|
||||
new_config["sample_size"] = 64
|
||||
unet._internal_dict = FrozenDict(new_config)
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
text_encoder=text_encoder,
|
||||
@@ -146,6 +176,8 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
def enable_attention_slicing(self, slice_size: Optional[Union[str, int]] = "auto"):
|
||||
r"""
|
||||
@@ -161,9 +193,14 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
if isinstance(self.unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.unet.config.attention_head_dim)
|
||||
|
||||
self.unet.set_attention_slice(slice_size)
|
||||
|
||||
def disable_attention_slicing(self):
|
||||
@@ -187,10 +224,15 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
|
||||
device = torch.device(f"cuda:{gpu_id}")
|
||||
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae, self.safety_checker]:
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae]:
|
||||
if cpu_offloaded_model is not None:
|
||||
cpu_offload(cpu_offloaded_model, device)
|
||||
|
||||
if self.safety_checker is not None:
|
||||
# TODO(Patrick) - there is currently a bug with cpu offload of nn.Parameter in accelerate
|
||||
# fix by only offloading self.safety_checker for now
|
||||
cpu_offload(self.safety_checker.vision_model, device)
|
||||
|
||||
@property
|
||||
def _execution_device(self):
|
||||
r"""
|
||||
@@ -391,11 +433,11 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
t_start = max(num_inference_steps - init_timestep + offset, 0)
|
||||
timesteps = self.scheduler.timesteps[t_start:]
|
||||
|
||||
return timesteps
|
||||
return timesteps, num_inference_steps - t_start
|
||||
|
||||
def prepare_latents(self, init_image, timestep, batch_size, num_images_per_prompt, dtype, device, generator=None):
|
||||
init_image = init_image.to(device=device, dtype=dtype)
|
||||
init_latent_dist = self.vae.encode(init_image).latent_dist
|
||||
def prepare_latents(self, image, timestep, batch_size, num_images_per_prompt, dtype, device, generator=None):
|
||||
image = image.to(device=device, dtype=dtype)
|
||||
init_latent_dist = self.vae.encode(image).latent_dist
|
||||
init_latents = init_latent_dist.sample(generator=generator)
|
||||
init_latents = 0.18215 * init_latents
|
||||
|
||||
@@ -403,16 +445,16 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
# expand init_latents for batch_size
|
||||
deprecation_message = (
|
||||
f"You have passed {batch_size} text prompts (`prompt`), but only {init_latents.shape[0]} initial"
|
||||
" images (`init_image`). Initial images are now duplicating to match the number of text prompts. Note"
|
||||
" images (`image`). Initial images are now duplicating to match the number of text prompts. Note"
|
||||
" that this behavior is deprecated and will be removed in a version 1.0.0. Please make sure to update"
|
||||
" your script to pass as many init images as text prompts to suppress this warning."
|
||||
" your script to pass as many initial images as text prompts to suppress this warning."
|
||||
)
|
||||
deprecate("len(prompt) != len(init_image)", "1.0.0", deprecation_message, standard_warn=False)
|
||||
deprecate("len(prompt) != len(image)", "1.0.0", deprecation_message, standard_warn=False)
|
||||
additional_image_per_prompt = batch_size // init_latents.shape[0]
|
||||
init_latents = torch.cat([init_latents] * additional_image_per_prompt * num_images_per_prompt, dim=0)
|
||||
elif batch_size > init_latents.shape[0] and batch_size % init_latents.shape[0] != 0:
|
||||
raise ValueError(
|
||||
f"Cannot duplicate `init_image` of batch size {init_latents.shape[0]} to {batch_size} text prompts."
|
||||
f"Cannot duplicate `image` of batch size {init_latents.shape[0]} to {batch_size} text prompts."
|
||||
)
|
||||
else:
|
||||
init_latents = torch.cat([init_latents] * num_images_per_prompt, dim=0)
|
||||
@@ -430,7 +472,7 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
init_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
strength: float = 0.8,
|
||||
num_inference_steps: Optional[int] = 50,
|
||||
guidance_scale: Optional[float] = 7.5,
|
||||
@@ -450,15 +492,15 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
init_image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process.
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
|
||||
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
|
||||
noise will be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
|
||||
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1. `image`
|
||||
will be used as a starting point, adding more noise to it the larger the `strength`. The number of
|
||||
denoising steps depends on the amount of noise initially added. When `strength` is 1, added noise will
|
||||
be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
expense of slower inference. This parameter will be modulated by `strength`.
|
||||
@@ -499,6 +541,10 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
message = "Please use `image` instead of `init_image`."
|
||||
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
|
||||
image = init_image or image
|
||||
|
||||
# 1. Check inputs
|
||||
self.check_inputs(prompt, strength, callback_steps)
|
||||
|
||||
@@ -516,42 +562,46 @@ class AltDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
)
|
||||
|
||||
# 4. Preprocess image
|
||||
if isinstance(init_image, PIL.Image.Image):
|
||||
init_image = preprocess(init_image)
|
||||
if isinstance(image, PIL.Image.Image):
|
||||
image = preprocess(image)
|
||||
|
||||
# 5. set timesteps
|
||||
self.scheduler.set_timesteps(num_inference_steps, device=device)
|
||||
timesteps = self.get_timesteps(num_inference_steps, strength, device)
|
||||
timesteps, num_inference_steps = self.get_timesteps(num_inference_steps, strength, device)
|
||||
latent_timestep = timesteps[:1].repeat(batch_size * num_images_per_prompt)
|
||||
|
||||
# 6. Prepare latent variables
|
||||
latents = self.prepare_latents(
|
||||
init_image, latent_timestep, batch_size, num_images_per_prompt, text_embeddings.dtype, device, generator
|
||||
image, latent_timestep, batch_size, num_images_per_prompt, text_embeddings.dtype, device, generator
|
||||
)
|
||||
|
||||
# 7. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
|
||||
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
|
||||
|
||||
# 8. Denoising loop
|
||||
for i, t in enumerate(self.progress_bar(timesteps)):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
for i, t in enumerate(timesteps):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
|
||||
# call the callback, if provided
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
# call the callback, if provided
|
||||
if (i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0:
|
||||
progress_bar.update()
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
|
||||
# 9. Post-processing
|
||||
image = self.decode_latents(latents)
|
||||
|
||||
@@ -70,14 +70,14 @@ class DDPMPipeline(DiffusionPipeline):
|
||||
generated images.
|
||||
"""
|
||||
message = (
|
||||
"Please make sure to instantiate your scheduler with `predict_epsilon` instead. E.g. `scheduler ="
|
||||
" DDPMScheduler.from_pretrained(<model_id>, predict_epsilon=True)`."
|
||||
"Please make sure to instantiate your scheduler with `prediction_type` instead. E.g. `scheduler ="
|
||||
" DDPMScheduler.from_pretrained(<model_id>, prediction_type='epsilon')`."
|
||||
)
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.10.0", message, take_from=kwargs)
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.11.0", message, take_from=kwargs)
|
||||
|
||||
if predict_epsilon is not None:
|
||||
new_config = dict(self.scheduler.config)
|
||||
new_config["predict_epsilon"] = predict_epsilon
|
||||
new_config["prediction_type"] = "epsilon" if predict_epsilon else "sample"
|
||||
self.scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if generator is not None and generator.device.type != self.device.type and self.device.type != "mps":
|
||||
@@ -114,9 +114,7 @@ class DDPMPipeline(DiffusionPipeline):
|
||||
model_output = self.unet(image, t).sample
|
||||
|
||||
# 2. compute previous image: x_t -> x_t-1
|
||||
image = self.scheduler.step(
|
||||
model_output, t, image, generator=generator, predict_epsilon=predict_epsilon
|
||||
).prev_sample
|
||||
image = self.scheduler.step(model_output, t, image, generator=generator).prev_sample
|
||||
|
||||
image = (image / 2 + 0.5).clamp(0, 1)
|
||||
image = image.cpu().permute(0, 2, 3, 1).numpy()
|
||||
|
||||
@@ -60,13 +60,14 @@ class LDMTextToImagePipeline(DiffusionPipeline):
|
||||
):
|
||||
super().__init__()
|
||||
self.register_modules(vqvae=vqvae, bert=bert, tokenizer=tokenizer, unet=unet, scheduler=scheduler)
|
||||
self.vae_scale_factor = 2 ** (len(self.vqvae.config.block_out_channels) - 1)
|
||||
|
||||
@torch.no_grad()
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
height: Optional[int] = 256,
|
||||
width: Optional[int] = 256,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: Optional[int] = 50,
|
||||
guidance_scale: Optional[float] = 1.0,
|
||||
eta: Optional[float] = 0.0,
|
||||
@@ -79,9 +80,9 @@ class LDMTextToImagePipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 256):
|
||||
height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 256):
|
||||
width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -106,6 +107,9 @@ class LDMTextToImagePipeline(DiffusionPipeline):
|
||||
`return_dict` is True, otherwise a `tuple. When returning a tuple, the first element is a list with the
|
||||
generated images.
|
||||
"""
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
if isinstance(prompt, str):
|
||||
batch_size = 1
|
||||
|
||||
@@ -17,7 +17,7 @@ from ...schedulers import (
|
||||
LMSDiscreteScheduler,
|
||||
PNDMScheduler,
|
||||
)
|
||||
from ...utils import PIL_INTERPOLATION
|
||||
from ...utils import PIL_INTERPOLATION, deprecate
|
||||
|
||||
|
||||
def preprocess(image):
|
||||
@@ -66,7 +66,7 @@ class LDMSuperResolutionPipeline(DiffusionPipeline):
|
||||
@torch.no_grad()
|
||||
def __call__(
|
||||
self,
|
||||
init_image: Union[torch.Tensor, PIL.Image.Image],
|
||||
image: Union[torch.Tensor, PIL.Image.Image],
|
||||
batch_size: Optional[int] = 1,
|
||||
num_inference_steps: Optional[int] = 100,
|
||||
eta: Optional[float] = 0.0,
|
||||
@@ -77,7 +77,7 @@ class LDMSuperResolutionPipeline(DiffusionPipeline):
|
||||
) -> Union[Tuple, ImagePipelineOutput]:
|
||||
r"""
|
||||
Args:
|
||||
init_image (`torch.Tensor` or `PIL.Image.Image`):
|
||||
image (`torch.Tensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process.
|
||||
batch_size (`int`, *optional*, defaults to 1):
|
||||
@@ -102,20 +102,21 @@ class LDMSuperResolutionPipeline(DiffusionPipeline):
|
||||
`return_dict` is True, otherwise a `tuple. When returning a tuple, the first element is a list with the
|
||||
generated images.
|
||||
"""
|
||||
message = "Please use `image` instead of `init_image`."
|
||||
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
|
||||
image = init_image or image
|
||||
|
||||
if isinstance(init_image, PIL.Image.Image):
|
||||
if isinstance(image, PIL.Image.Image):
|
||||
batch_size = 1
|
||||
elif isinstance(init_image, torch.Tensor):
|
||||
batch_size = init_image.shape[0]
|
||||
elif isinstance(image, torch.Tensor):
|
||||
batch_size = image.shape[0]
|
||||
else:
|
||||
raise ValueError(
|
||||
f"`init_image` has to be of type `PIL.Image.Image` or `torch.Tensor` but is {type(init_image)}"
|
||||
)
|
||||
raise ValueError(f"`image` has to be of type `PIL.Image.Image` or `torch.Tensor` but is {type(image)}")
|
||||
|
||||
if isinstance(init_image, PIL.Image.Image):
|
||||
init_image = preprocess(init_image)
|
||||
if isinstance(image, PIL.Image.Image):
|
||||
image = preprocess(image)
|
||||
|
||||
height, width = init_image.shape[-2:]
|
||||
height, width = image.shape[-2:]
|
||||
|
||||
# in_channels should be 6: 3 for latents, 3 for low resolution image
|
||||
latents_shape = (batch_size, self.unet.in_channels // 2, height, width)
|
||||
@@ -128,7 +129,7 @@ class LDMSuperResolutionPipeline(DiffusionPipeline):
|
||||
else:
|
||||
latents = torch.randn(latents_shape, generator=generator, device=self.device, dtype=latents_dtype)
|
||||
|
||||
init_image = init_image.to(device=self.device, dtype=latents_dtype)
|
||||
image = image.to(device=self.device, dtype=latents_dtype)
|
||||
|
||||
# set timesteps and move to the correct device
|
||||
self.scheduler.set_timesteps(num_inference_steps, device=self.device)
|
||||
@@ -148,7 +149,7 @@ class LDMSuperResolutionPipeline(DiffusionPipeline):
|
||||
|
||||
for t in self.progress_bar(timesteps_tensor):
|
||||
# concat latents and low resolution image in the channel dimension.
|
||||
latents_input = torch.cat([latents, init_image], dim=1)
|
||||
latents_input = torch.cat([latents, image], dim=1)
|
||||
latents_input = self.scheduler.scale_model_input(latents_input, t)
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latents_input, t).sample
|
||||
|
||||
@@ -138,7 +138,7 @@ prompt = "An astronaut riding an elephant"
|
||||
image = pipe(
|
||||
prompt=prompt,
|
||||
source_prompt=source_prompt,
|
||||
init_image=init_image,
|
||||
image=init_image,
|
||||
num_inference_steps=100,
|
||||
eta=0.1,
|
||||
strength=0.8,
|
||||
@@ -164,7 +164,7 @@ torch.manual_seed(0)
|
||||
image = pipe(
|
||||
prompt=prompt,
|
||||
source_prompt=source_prompt,
|
||||
init_image=init_image,
|
||||
image=init_image,
|
||||
num_inference_steps=100,
|
||||
eta=0.1,
|
||||
strength=0.85,
|
||||
|
||||
@@ -40,6 +40,7 @@ if is_transformers_available() and is_torch_available():
|
||||
from .pipeline_stable_diffusion_img2img import StableDiffusionImg2ImgPipeline
|
||||
from .pipeline_stable_diffusion_inpaint import StableDiffusionInpaintPipeline
|
||||
from .pipeline_stable_diffusion_inpaint_legacy import StableDiffusionInpaintPipelineLegacy
|
||||
from .pipeline_stable_diffusion_upscale import StableDiffusionUpscalePipeline
|
||||
from .safety_checker import StableDiffusionSafetyChecker
|
||||
|
||||
if is_transformers_available() and is_torch_available() and is_transformers_version(">=", "4.25.0.dev0"):
|
||||
@@ -62,15 +63,14 @@ if is_transformers_available() and is_flax_available():
|
||||
Output class for Stable Diffusion pipelines.
|
||||
|
||||
Args:
|
||||
images (`List[PIL.Image.Image]` or `np.ndarray`)
|
||||
List of denoised PIL images of length `batch_size` or numpy array of shape `(batch_size, height, width,
|
||||
num_channels)`. PIL images or numpy array present the denoised images of the diffusion pipeline.
|
||||
images (`np.ndarray`)
|
||||
Array of shape `(batch_size, height, width, num_channels)` with images from the diffusion pipeline.
|
||||
nsfw_content_detected (`List[bool]`)
|
||||
List of flags denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content.
|
||||
"""
|
||||
|
||||
images: Union[List[PIL.Image.Image], np.ndarray]
|
||||
images: np.ndarray
|
||||
nsfw_content_detected: List[bool]
|
||||
|
||||
from ...schedulers.scheduling_pndm_flax import PNDMSchedulerState
|
||||
|
||||
@@ -20,6 +20,7 @@ import torch
|
||||
|
||||
import PIL
|
||||
from diffusers.utils import is_accelerate_available
|
||||
from packaging import version
|
||||
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
|
||||
|
||||
from ...configuration_utils import FrozenDict
|
||||
@@ -132,6 +133,7 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
feature_extractor ([`CLIPFeatureExtractor`]):
|
||||
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
|
||||
"""
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
@@ -142,6 +144,7 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
scheduler: DDIMScheduler,
|
||||
safety_checker: StableDiffusionSafetyChecker,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -159,7 +162,7 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
new_config["steps_offset"] = 1
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
@@ -169,6 +172,32 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
|
||||
version.parse(unet.config._diffusers_version).base_version
|
||||
) < version.parse("0.9.0.dev0")
|
||||
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
|
||||
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
|
||||
deprecation_message = (
|
||||
"The configuration file of the unet has set the default `sample_size` to smaller than"
|
||||
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
|
||||
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
|
||||
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
|
||||
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
|
||||
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
|
||||
" in the config might lead to incorrect results in future versions. If you have downloaded this"
|
||||
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
|
||||
" the `unet/config.json` file"
|
||||
)
|
||||
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
|
||||
new_config = dict(unet.config)
|
||||
new_config["sample_size"] = 64
|
||||
unet._internal_dict = FrozenDict(new_config)
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
text_encoder=text_encoder,
|
||||
@@ -178,6 +207,7 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_attention_slicing
|
||||
def enable_attention_slicing(self, slice_size: Optional[Union[str, int]] = "auto"):
|
||||
@@ -194,9 +224,14 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
if isinstance(self.unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.unet.config.attention_head_dim)
|
||||
|
||||
self.unet.set_attention_slice(slice_size)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_attention_slicing
|
||||
@@ -222,10 +257,15 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
|
||||
device = torch.device(f"cuda:{gpu_id}")
|
||||
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae, self.safety_checker]:
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae]:
|
||||
if cpu_offloaded_model is not None:
|
||||
cpu_offload(cpu_offloaded_model, device)
|
||||
|
||||
if self.safety_checker is not None:
|
||||
# TODO(Patrick) - there is currently a bug with cpu offload of nn.Parameter in accelerate
|
||||
# fix by only offloading self.safety_checker for now
|
||||
cpu_offload(self.safety_checker.vision_model, device)
|
||||
|
||||
@property
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline._execution_device
|
||||
def _execution_device(self):
|
||||
@@ -435,11 +475,11 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
t_start = max(num_inference_steps - init_timestep + offset, 0)
|
||||
timesteps = self.scheduler.timesteps[t_start:]
|
||||
|
||||
return timesteps
|
||||
return timesteps, num_inference_steps - t_start
|
||||
|
||||
def prepare_latents(self, init_image, timestep, batch_size, num_images_per_prompt, dtype, device, generator=None):
|
||||
init_image = init_image.to(device=device, dtype=dtype)
|
||||
init_latent_dist = self.vae.encode(init_image).latent_dist
|
||||
def prepare_latents(self, image, timestep, batch_size, num_images_per_prompt, dtype, device, generator=None):
|
||||
image = image.to(device=device, dtype=dtype)
|
||||
init_latent_dist = self.vae.encode(image).latent_dist
|
||||
init_latents = init_latent_dist.sample(generator=generator)
|
||||
init_latents = 0.18215 * init_latents
|
||||
|
||||
@@ -447,16 +487,16 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
# expand init_latents for batch_size
|
||||
deprecation_message = (
|
||||
f"You have passed {batch_size} text prompts (`prompt`), but only {init_latents.shape[0]} initial"
|
||||
" images (`init_image`). Initial images are now duplicating to match the number of text prompts. Note"
|
||||
" images (`image`). Initial images are now duplicating to match the number of text prompts. Note"
|
||||
" that this behavior is deprecated and will be removed in a version 1.0.0. Please make sure to update"
|
||||
" your script to pass as many init images as text prompts to suppress this warning."
|
||||
" your script to pass as many initial images as text prompts to suppress this warning."
|
||||
)
|
||||
deprecate("len(prompt) != len(init_image)", "1.0.0", deprecation_message, standard_warn=False)
|
||||
deprecate("len(prompt) != len(image)", "1.0.0", deprecation_message, standard_warn=False)
|
||||
additional_image_per_prompt = batch_size // init_latents.shape[0]
|
||||
init_latents = torch.cat([init_latents] * additional_image_per_prompt * num_images_per_prompt, dim=0)
|
||||
elif batch_size > init_latents.shape[0] and batch_size % init_latents.shape[0] != 0:
|
||||
raise ValueError(
|
||||
f"Cannot duplicate `init_image` of batch size {init_latents.shape[0]} to {batch_size} text prompts."
|
||||
f"Cannot duplicate `image` of batch size {init_latents.shape[0]} to {batch_size} text prompts."
|
||||
)
|
||||
else:
|
||||
init_latents = torch.cat([init_latents] * num_images_per_prompt, dim=0)
|
||||
@@ -476,7 +516,7 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
source_prompt: Union[str, List[str]],
|
||||
init_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
strength: float = 0.8,
|
||||
num_inference_steps: Optional[int] = 50,
|
||||
guidance_scale: Optional[float] = 7.5,
|
||||
@@ -496,15 +536,15 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
init_image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process.
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
|
||||
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
|
||||
noise will be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
|
||||
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1. `image`
|
||||
will be used as a starting point, adding more noise to it the larger the `strength`. The number of
|
||||
denoising steps depends on the amount of noise initially added. When `strength` is 1, added noise will
|
||||
be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
expense of slower inference. This parameter will be modulated by `strength`.
|
||||
@@ -545,6 +585,10 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
message = "Please use `image` instead of `init_image`."
|
||||
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
|
||||
image = init_image or image
|
||||
|
||||
# 1. Check inputs
|
||||
self.check_inputs(prompt, strength, callback_steps)
|
||||
|
||||
@@ -563,17 +607,17 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
)
|
||||
|
||||
# 4. Preprocess image
|
||||
if isinstance(init_image, PIL.Image.Image):
|
||||
init_image = preprocess(init_image)
|
||||
if isinstance(image, PIL.Image.Image):
|
||||
image = preprocess(image)
|
||||
|
||||
# 5. Prepare timesteps
|
||||
self.scheduler.set_timesteps(num_inference_steps, device=device)
|
||||
timesteps = self.get_timesteps(num_inference_steps, strength, device)
|
||||
timesteps, num_inference_steps = self.get_timesteps(num_inference_steps, strength, device)
|
||||
latent_timestep = timesteps[:1].repeat(batch_size * num_images_per_prompt)
|
||||
|
||||
# 6. Prepare latent variables
|
||||
latents, clean_latents = self.prepare_latents(
|
||||
init_image, latent_timestep, batch_size, num_images_per_prompt, text_embeddings.dtype, device, generator
|
||||
image, latent_timestep, batch_size, num_images_per_prompt, text_embeddings.dtype, device, generator
|
||||
)
|
||||
source_latents = latents
|
||||
|
||||
@@ -582,66 +626,70 @@ class CycleDiffusionPipeline(DiffusionPipeline):
|
||||
generator = extra_step_kwargs.pop("generator", None)
|
||||
|
||||
# 8. Denoising loop
|
||||
for i, t in enumerate(self.progress_bar(timesteps)):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2)
|
||||
source_latent_model_input = torch.cat([source_latents] * 2)
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
source_latent_model_input = self.scheduler.scale_model_input(source_latent_model_input, t)
|
||||
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
for i, t in enumerate(timesteps):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2)
|
||||
source_latent_model_input = torch.cat([source_latents] * 2)
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
source_latent_model_input = self.scheduler.scale_model_input(source_latent_model_input, t)
|
||||
|
||||
# predict the noise residual
|
||||
concat_latent_model_input = torch.stack(
|
||||
[
|
||||
source_latent_model_input[0],
|
||||
latent_model_input[0],
|
||||
source_latent_model_input[1],
|
||||
latent_model_input[1],
|
||||
],
|
||||
dim=0,
|
||||
)
|
||||
concat_text_embeddings = torch.stack(
|
||||
[
|
||||
source_text_embeddings[0],
|
||||
text_embeddings[0],
|
||||
source_text_embeddings[1],
|
||||
text_embeddings[1],
|
||||
],
|
||||
dim=0,
|
||||
)
|
||||
concat_noise_pred = self.unet(
|
||||
concat_latent_model_input, t, encoder_hidden_states=concat_text_embeddings
|
||||
).sample
|
||||
# predict the noise residual
|
||||
concat_latent_model_input = torch.stack(
|
||||
[
|
||||
source_latent_model_input[0],
|
||||
latent_model_input[0],
|
||||
source_latent_model_input[1],
|
||||
latent_model_input[1],
|
||||
],
|
||||
dim=0,
|
||||
)
|
||||
concat_text_embeddings = torch.stack(
|
||||
[
|
||||
source_text_embeddings[0],
|
||||
text_embeddings[0],
|
||||
source_text_embeddings[1],
|
||||
text_embeddings[1],
|
||||
],
|
||||
dim=0,
|
||||
)
|
||||
concat_noise_pred = self.unet(
|
||||
concat_latent_model_input, t, encoder_hidden_states=concat_text_embeddings
|
||||
).sample
|
||||
|
||||
# perform guidance
|
||||
(
|
||||
source_noise_pred_uncond,
|
||||
noise_pred_uncond,
|
||||
source_noise_pred_text,
|
||||
noise_pred_text,
|
||||
) = concat_noise_pred.chunk(4, dim=0)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
source_noise_pred = source_noise_pred_uncond + source_guidance_scale * (
|
||||
source_noise_pred_text - source_noise_pred_uncond
|
||||
)
|
||||
# perform guidance
|
||||
(
|
||||
source_noise_pred_uncond,
|
||||
noise_pred_uncond,
|
||||
source_noise_pred_text,
|
||||
noise_pred_text,
|
||||
) = concat_noise_pred.chunk(4, dim=0)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
source_noise_pred = source_noise_pred_uncond + source_guidance_scale * (
|
||||
source_noise_pred_text - source_noise_pred_uncond
|
||||
)
|
||||
|
||||
# Sample source_latents from the posterior distribution.
|
||||
prev_source_latents = posterior_sample(
|
||||
self.scheduler, source_latents, t, clean_latents, generator=generator, **extra_step_kwargs
|
||||
)
|
||||
# Compute noise.
|
||||
noise = compute_noise(
|
||||
self.scheduler, prev_source_latents, source_latents, t, source_noise_pred, **extra_step_kwargs
|
||||
)
|
||||
source_latents = prev_source_latents
|
||||
# Sample source_latents from the posterior distribution.
|
||||
prev_source_latents = posterior_sample(
|
||||
self.scheduler, source_latents, t, clean_latents, generator=generator, **extra_step_kwargs
|
||||
)
|
||||
# Compute noise.
|
||||
noise = compute_noise(
|
||||
self.scheduler, prev_source_latents, source_latents, t, source_noise_pred, **extra_step_kwargs
|
||||
)
|
||||
source_latents = prev_source_latents
|
||||
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(
|
||||
noise_pred, t, latents, variance_noise=noise, **extra_step_kwargs
|
||||
).prev_sample
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(
|
||||
noise_pred, t, latents, variance_noise=noise, **extra_step_kwargs
|
||||
).prev_sample
|
||||
|
||||
# call the callback, if provided
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
# call the callback, if provided
|
||||
if (i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0:
|
||||
progress_bar.update()
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
|
||||
# 9. Post-processing
|
||||
image = self.decode_latents(latents)
|
||||
|
||||
@@ -23,6 +23,7 @@ import jax.numpy as jnp
|
||||
from flax.core.frozen_dict import FrozenDict
|
||||
from flax.jax_utils import unreplicate
|
||||
from flax.training.common_utils import shard
|
||||
from packaging import version
|
||||
from PIL import Image
|
||||
from transformers import CLIPFeatureExtractor, CLIPTokenizer, FlaxCLIPTextModel
|
||||
|
||||
@@ -34,7 +35,7 @@ from ...schedulers import (
|
||||
FlaxLMSDiscreteScheduler,
|
||||
FlaxPNDMScheduler,
|
||||
)
|
||||
from ...utils import logging
|
||||
from ...utils import deprecate, logging
|
||||
from . import FlaxStableDiffusionPipelineOutput
|
||||
from .safety_checker_flax import FlaxStableDiffusionSafetyChecker
|
||||
|
||||
@@ -97,6 +98,27 @@ class FlaxStableDiffusionPipeline(FlaxDiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
|
||||
version.parse(unet.config._diffusers_version).base_version
|
||||
) < version.parse("0.9.0.dev0")
|
||||
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
|
||||
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
|
||||
deprecation_message = (
|
||||
"The configuration file of the unet has set the default `sample_size` to smaller than"
|
||||
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
|
||||
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
|
||||
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
|
||||
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
|
||||
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
|
||||
" in the config might lead to incorrect results in future versions. If you have downloaded this"
|
||||
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
|
||||
" the `unet/config.json` file"
|
||||
)
|
||||
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
|
||||
new_config = dict(unet.config)
|
||||
new_config["sample_size"] = 64
|
||||
unet._internal_dict = FrozenDict(new_config)
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
text_encoder=text_encoder,
|
||||
@@ -106,6 +128,7 @@ class FlaxStableDiffusionPipeline(FlaxDiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
|
||||
def prepare_inputs(self, prompt: Union[str, List[str]]):
|
||||
if not isinstance(prompt, (str, list)):
|
||||
@@ -160,13 +183,17 @@ class FlaxStableDiffusionPipeline(FlaxDiffusionPipeline):
|
||||
params: Union[Dict, FrozenDict],
|
||||
prng_seed: jax.random.PRNGKey,
|
||||
num_inference_steps: int = 50,
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
guidance_scale: float = 7.5,
|
||||
latents: Optional[jnp.array] = None,
|
||||
debug: bool = False,
|
||||
neg_prompt_ids: jnp.array = None,
|
||||
):
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
if height % 8 != 0 or width % 8 != 0:
|
||||
raise ValueError(f"`height` and `width` have to be divisible by 8 but are {height} and {width}.")
|
||||
|
||||
@@ -188,7 +215,12 @@ class FlaxStableDiffusionPipeline(FlaxDiffusionPipeline):
|
||||
uncond_embeddings = self.text_encoder(uncond_input, params=params["text_encoder"])[0]
|
||||
context = jnp.concatenate([uncond_embeddings, text_embeddings])
|
||||
|
||||
latents_shape = (batch_size, self.unet.in_channels, height // 8, width // 8)
|
||||
latents_shape = (
|
||||
batch_size,
|
||||
self.unet.in_channels,
|
||||
height // self.vae_scale_factor,
|
||||
width // self.vae_scale_factor,
|
||||
)
|
||||
if latents is None:
|
||||
latents = jax.random.normal(prng_seed, shape=latents_shape, dtype=jnp.float32)
|
||||
else:
|
||||
@@ -249,15 +281,14 @@ class FlaxStableDiffusionPipeline(FlaxDiffusionPipeline):
|
||||
params: Union[Dict, FrozenDict],
|
||||
prng_seed: jax.random.PRNGKey,
|
||||
num_inference_steps: int = 50,
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
guidance_scale: float = 7.5,
|
||||
latents: jnp.array = None,
|
||||
return_dict: bool = True,
|
||||
jit: bool = False,
|
||||
debug: bool = False,
|
||||
neg_prompt_ids: jnp.array = None,
|
||||
**kwargs,
|
||||
):
|
||||
r"""
|
||||
Function invoked when calling the pipeline for generation.
|
||||
@@ -265,9 +296,9 @@ class FlaxStableDiffusionPipeline(FlaxDiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -285,9 +316,6 @@ class FlaxStableDiffusionPipeline(FlaxDiffusionPipeline):
|
||||
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
|
||||
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
|
||||
tensor will ge generated by sampling using the supplied random `generator`.
|
||||
output_type (`str`, *optional*, defaults to `"pil"`):
|
||||
The output format of the generate image. Choose between
|
||||
[PIL](https://pillow.readthedocs.io/en/stable/): `PIL.Image.Image` or `np.array`.
|
||||
jit (`bool`, defaults to `False`):
|
||||
Whether to run `pmap` versions of the generation and safety scoring functions. NOTE: This argument
|
||||
exists because `__call__` is not yet end-to-end pmap-able. It will be removed in a future release.
|
||||
@@ -302,6 +330,10 @@ class FlaxStableDiffusionPipeline(FlaxDiffusionPipeline):
|
||||
element is a list of `bool`s denoting whether the corresponding generated image likely represents
|
||||
"not-safe-for-work" (nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
if jit:
|
||||
images = _p_generate(
|
||||
self,
|
||||
@@ -347,6 +379,7 @@ class FlaxStableDiffusionPipeline(FlaxDiffusionPipeline):
|
||||
|
||||
images = images.reshape(num_devices, batch_size, height, width, 3)
|
||||
else:
|
||||
images = np.asarray(images)
|
||||
has_nsfw_concept = False
|
||||
|
||||
if not return_dict:
|
||||
|
||||
@@ -41,6 +41,8 @@ class OnnxStableDiffusionPipeline(DiffusionPipeline):
|
||||
safety_checker: OnnxRuntimeModel
|
||||
feature_extractor: CLIPFeatureExtractor
|
||||
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
vae_encoder: OnnxRuntimeModel,
|
||||
@@ -51,6 +53,7 @@ class OnnxStableDiffusionPipeline(DiffusionPipeline):
|
||||
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
|
||||
safety_checker: OnnxRuntimeModel,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -81,6 +84,22 @@ class OnnxStableDiffusionPipeline(DiffusionPipeline):
|
||||
new_config["clip_sample"] = False
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
" results in services or applications open to the public. Both the diffusers team and Hugging Face"
|
||||
" strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling"
|
||||
" it only for use-cases that involve analyzing network behavior or auditing its results. For more"
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
self.register_modules(
|
||||
vae_encoder=vae_encoder,
|
||||
vae_decoder=vae_decoder,
|
||||
@@ -91,6 +110,7 @@ class OnnxStableDiffusionPipeline(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
def _encode_prompt(self, prompt, num_images_per_prompt, do_classifier_free_guidance, negative_prompt):
|
||||
r"""
|
||||
@@ -185,7 +205,6 @@ class OnnxStableDiffusionPipeline(DiffusionPipeline):
|
||||
return_dict: bool = True,
|
||||
callback: Optional[Callable[[int, int, np.ndarray], None]] = None,
|
||||
callback_steps: Optional[int] = 1,
|
||||
**kwargs,
|
||||
):
|
||||
if isinstance(prompt, str):
|
||||
batch_size = 1
|
||||
|
||||
@@ -77,6 +77,8 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
safety_checker: OnnxRuntimeModel
|
||||
feature_extractor: CLIPFeatureExtractor
|
||||
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
vae_encoder: OnnxRuntimeModel,
|
||||
@@ -87,6 +89,7 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
|
||||
safety_checker: OnnxRuntimeModel,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -117,7 +120,7 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
new_config["clip_sample"] = False
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
@@ -127,6 +130,12 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
self.register_modules(
|
||||
vae_encoder=vae_encoder,
|
||||
vae_decoder=vae_decoder,
|
||||
@@ -137,6 +146,7 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_onnx_stable_diffusion.OnnxStableDiffusionPipeline._encode_prompt
|
||||
def _encode_prompt(self, prompt, num_images_per_prompt, do_classifier_free_guidance, negative_prompt):
|
||||
@@ -219,7 +229,7 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
init_image: Union[np.ndarray, PIL.Image.Image],
|
||||
image: Union[np.ndarray, PIL.Image.Image],
|
||||
strength: float = 0.8,
|
||||
num_inference_steps: Optional[int] = 50,
|
||||
guidance_scale: Optional[float] = 7.5,
|
||||
@@ -239,15 +249,15 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
init_image (`np.ndarray` or `PIL.Image.Image`):
|
||||
image (`np.ndarray` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process.
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
|
||||
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
|
||||
noise will be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
|
||||
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1. `image`
|
||||
will be used as a starting point, adding more noise to it the larger the `strength`. The number of
|
||||
denoising steps depends on the amount of noise initially added. When `strength` is 1, added noise will
|
||||
be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
expense of slower inference. This parameter will be modulated by `strength`.
|
||||
@@ -287,6 +297,10 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
message = "Please use `image` instead of `init_image`."
|
||||
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
|
||||
image = init_image or image
|
||||
|
||||
if isinstance(prompt, str):
|
||||
batch_size = 1
|
||||
elif isinstance(prompt, list):
|
||||
@@ -311,8 +325,8 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
# set timesteps
|
||||
self.scheduler.set_timesteps(num_inference_steps)
|
||||
|
||||
if isinstance(init_image, PIL.Image.Image):
|
||||
init_image = preprocess(init_image)
|
||||
if isinstance(image, PIL.Image.Image):
|
||||
image = preprocess(image)
|
||||
|
||||
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
|
||||
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
|
||||
@@ -324,9 +338,9 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
)
|
||||
|
||||
latents_dtype = text_embeddings.dtype
|
||||
init_image = init_image.astype(latents_dtype)
|
||||
image = image.astype(latents_dtype)
|
||||
# encode the init image into latents and scale the latents
|
||||
init_latents = self.vae_encoder(sample=init_image)[0]
|
||||
init_latents = self.vae_encoder(sample=image)[0]
|
||||
init_latents = 0.18215 * init_latents
|
||||
|
||||
if isinstance(prompt, str):
|
||||
@@ -335,16 +349,16 @@ class OnnxStableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
# expand init_latents for batch_size
|
||||
deprecation_message = (
|
||||
f"You have passed {len(prompt)} text prompts (`prompt`), but only {init_latents.shape[0]} initial"
|
||||
" images (`init_image`). Initial images are now duplicating to match the number of text prompts. Note"
|
||||
" images (`image`). Initial images are now duplicating to match the number of text prompts. Note"
|
||||
" that this behavior is deprecated and will be removed in a version 1.0.0. Please make sure to update"
|
||||
" your script to pass as many init images as text prompts to suppress this warning."
|
||||
" your script to pass as many initial images as text prompts to suppress this warning."
|
||||
)
|
||||
deprecate("len(prompt) != len(init_image)", "1.0.0", deprecation_message, standard_warn=False)
|
||||
deprecate("len(prompt) != len(image)", "1.0.0", deprecation_message, standard_warn=False)
|
||||
additional_image_per_prompt = len(prompt) // init_latents.shape[0]
|
||||
init_latents = np.concatenate([init_latents] * additional_image_per_prompt * num_images_per_prompt, axis=0)
|
||||
elif len(prompt) > init_latents.shape[0] and len(prompt) % init_latents.shape[0] != 0:
|
||||
raise ValueError(
|
||||
f"Cannot duplicate `init_image` of batch size {init_latents.shape[0]} to {len(prompt)} text prompts."
|
||||
f"Cannot duplicate `image` of batch size {init_latents.shape[0]} to {len(prompt)} text prompts."
|
||||
)
|
||||
else:
|
||||
init_latents = np.concatenate([init_latents] * num_images_per_prompt, axis=0)
|
||||
|
||||
@@ -90,6 +90,8 @@ class OnnxStableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
safety_checker: OnnxRuntimeModel
|
||||
feature_extractor: CLIPFeatureExtractor
|
||||
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
vae_encoder: OnnxRuntimeModel,
|
||||
@@ -100,6 +102,7 @@ class OnnxStableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
|
||||
safety_checker: OnnxRuntimeModel,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
logger.info("`OnnxStableDiffusionInpaintPipeline` is experimental and will very likely change in the future.")
|
||||
@@ -131,7 +134,7 @@ class OnnxStableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
new_config["clip_sample"] = False
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
@@ -141,6 +144,12 @@ class OnnxStableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
self.register_modules(
|
||||
vae_encoder=vae_encoder,
|
||||
vae_decoder=vae_decoder,
|
||||
@@ -151,6 +160,7 @@ class OnnxStableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_onnx_stable_diffusion.OnnxStableDiffusionPipeline._encode_prompt
|
||||
def _encode_prompt(self, prompt, num_images_per_prompt, do_classifier_free_guidance, negative_prompt):
|
||||
@@ -236,8 +246,8 @@ class OnnxStableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
prompt: Union[str, List[str]],
|
||||
image: PIL.Image.Image,
|
||||
mask_image: PIL.Image.Image,
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = 512,
|
||||
width: Optional[int] = 512,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -249,7 +259,6 @@ class OnnxStableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
return_dict: bool = True,
|
||||
callback: Optional[Callable[[int, int, np.ndarray], None]] = None,
|
||||
callback_steps: Optional[int] = 1,
|
||||
**kwargs,
|
||||
):
|
||||
r"""
|
||||
Function invoked when calling the pipeline for generation.
|
||||
@@ -312,6 +321,7 @@ class OnnxStableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
|
||||
if isinstance(prompt, str):
|
||||
batch_size = 1
|
||||
elif isinstance(prompt, list):
|
||||
|
||||
@@ -27,11 +27,11 @@ def preprocess(image):
|
||||
return 2.0 * image - 1.0
|
||||
|
||||
|
||||
def preprocess_mask(mask):
|
||||
def preprocess_mask(mask, scale_factor=8):
|
||||
mask = mask.convert("L")
|
||||
w, h = mask.size
|
||||
w, h = map(lambda x: x - x % 32, (w, h)) # resize to integer multiple of 32
|
||||
mask = mask.resize((w // 8, h // 8), resample=PIL.Image.NEAREST)
|
||||
mask = mask.resize((w // scale_factor, h // scale_factor), resample=PIL.Image.NEAREST)
|
||||
mask = np.array(mask).astype(np.float32) / 255.0
|
||||
mask = np.tile(mask, (4, 1, 1))
|
||||
mask = mask[None].transpose(0, 1, 2, 3) # what does this step do?
|
||||
@@ -67,6 +67,8 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
feature_extractor ([`CLIPFeatureExtractor`]):
|
||||
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
|
||||
"""
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
vae_encoder: OnnxRuntimeModel
|
||||
vae_decoder: OnnxRuntimeModel
|
||||
text_encoder: OnnxRuntimeModel
|
||||
@@ -86,6 +88,7 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
|
||||
safety_checker: OnnxRuntimeModel,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -116,7 +119,7 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
new_config["clip_sample"] = False
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
@@ -126,6 +129,12 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
self.register_modules(
|
||||
vae_encoder=vae_encoder,
|
||||
vae_decoder=vae_decoder,
|
||||
@@ -136,6 +145,7 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_onnx_stable_diffusion.OnnxStableDiffusionPipeline._encode_prompt
|
||||
def _encode_prompt(self, prompt, num_images_per_prompt, do_classifier_free_guidance, negative_prompt):
|
||||
@@ -218,7 +228,7 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
init_image: Union[np.ndarray, PIL.Image.Image],
|
||||
image: Union[np.ndarray, PIL.Image.Image],
|
||||
mask_image: Union[np.ndarray, PIL.Image.Image],
|
||||
strength: float = 0.8,
|
||||
num_inference_steps: Optional[int] = 50,
|
||||
@@ -239,20 +249,20 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
init_image (`nd.ndarray` or `PIL.Image.Image`):
|
||||
image (`nd.ndarray` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process. This is the image whose masked region will be inpainted.
|
||||
mask_image (`nd.ndarray` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, to mask `init_image`. White pixels in the mask will be
|
||||
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be
|
||||
replaced by noise and therefore repainted, while black pixels will be preserved. If `mask_image` is a
|
||||
PIL image, it will be converted to a single channel (luminance) before use. If it's a tensor, it should
|
||||
contain one color channel (L) instead of 3, so the expected shape would be `(B, H, W, 1)`.uu
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
|
||||
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
|
||||
noise will be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
|
||||
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1. `image`
|
||||
will be used as a starting point, adding more noise to it the larger the `strength`. The number of
|
||||
denoising steps depends on the amount of noise initially added. When `strength` is 1, added noise will
|
||||
be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
expense of slower inference. This parameter will be modulated by `strength`.
|
||||
@@ -292,6 +302,10 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
message = "Please use `image` instead of `init_image`."
|
||||
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
|
||||
image = init_image or image
|
||||
|
||||
if isinstance(prompt, str):
|
||||
batch_size = 1
|
||||
elif isinstance(prompt, list):
|
||||
@@ -316,8 +330,8 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
# set timesteps
|
||||
self.scheduler.set_timesteps(num_inference_steps)
|
||||
|
||||
if isinstance(init_image, PIL.Image.Image):
|
||||
init_image = preprocess(init_image)
|
||||
if isinstance(image, PIL.Image.Image):
|
||||
image = preprocess(image)
|
||||
|
||||
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
|
||||
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
|
||||
@@ -329,10 +343,10 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
)
|
||||
|
||||
latents_dtype = text_embeddings.dtype
|
||||
init_image = init_image.astype(latents_dtype)
|
||||
image = image.astype(latents_dtype)
|
||||
|
||||
# encode the init image into latents and scale the latents
|
||||
init_latents = self.vae_encoder(sample=init_image)[0]
|
||||
init_latents = self.vae_encoder(sample=image)[0]
|
||||
init_latents = 0.18215 * init_latents
|
||||
|
||||
# Expand init_latents for batch_size and num_images_per_prompt
|
||||
@@ -341,13 +355,13 @@ class OnnxStableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
|
||||
# preprocess mask
|
||||
if not isinstance(mask_image, np.ndarray):
|
||||
mask_image = preprocess_mask(mask_image)
|
||||
mask_image = preprocess_mask(mask_image, 8)
|
||||
mask_image = mask_image.astype(latents_dtype)
|
||||
mask = np.concatenate([mask_image] * num_images_per_prompt, axis=0)
|
||||
|
||||
# check sizes
|
||||
if not mask.shape == init_latents.shape:
|
||||
raise ValueError("The mask and init_image should be the same size!")
|
||||
raise ValueError("The mask and image should be the same size!")
|
||||
|
||||
# get the original timestep using init_timestep
|
||||
offset = self.scheduler.config.get("steps_offset", 0)
|
||||
|
||||
@@ -18,6 +18,7 @@ from typing import Callable, List, Optional, Union
|
||||
import torch
|
||||
|
||||
from diffusers.utils import is_accelerate_available
|
||||
from packaging import version
|
||||
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
|
||||
|
||||
from ...configuration_utils import FrozenDict
|
||||
@@ -66,6 +67,7 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
feature_extractor ([`CLIPFeatureExtractor`]):
|
||||
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
|
||||
"""
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
@@ -83,6 +85,7 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
],
|
||||
safety_checker: StableDiffusionSafetyChecker,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -113,7 +116,7 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
new_config["clip_sample"] = False
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
@@ -123,6 +126,33 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
|
||||
version.parse(unet.config._diffusers_version).base_version
|
||||
) < version.parse("0.9.0.dev0")
|
||||
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
|
||||
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
|
||||
deprecation_message = (
|
||||
"The configuration file of the unet has set the default `sample_size` to smaller than"
|
||||
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
|
||||
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
|
||||
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
|
||||
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
|
||||
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
|
||||
" in the config might lead to incorrect results in future versions. If you have downloaded this"
|
||||
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
|
||||
" the `unet/config.json` file"
|
||||
)
|
||||
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
|
||||
new_config = dict(unet.config)
|
||||
new_config["sample_size"] = 64
|
||||
unet._internal_dict = FrozenDict(new_config)
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
text_encoder=text_encoder,
|
||||
@@ -132,6 +162,8 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
def enable_xformers_memory_efficient_attention(self):
|
||||
r"""
|
||||
@@ -165,9 +197,14 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
if isinstance(self.unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.unet.config.attention_head_dim)
|
||||
|
||||
self.unet.set_attention_slice(slice_size)
|
||||
|
||||
def disable_attention_slicing(self):
|
||||
@@ -178,6 +215,22 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
# set slice_size = `None` to disable `attention slicing`
|
||||
self.enable_attention_slicing(None)
|
||||
|
||||
def enable_vae_slicing(self):
|
||||
r"""
|
||||
Enable sliced VAE decoding.
|
||||
|
||||
When this option is enabled, the VAE will split the input tensor in slices to compute decoding in several
|
||||
steps. This is useful to save some memory and allow larger batch sizes.
|
||||
"""
|
||||
self.vae.enable_slicing()
|
||||
|
||||
def disable_vae_slicing(self):
|
||||
r"""
|
||||
Disable sliced VAE decoding. If `enable_vae_slicing` was previously invoked, this method will go back to
|
||||
computing decoding in one step.
|
||||
"""
|
||||
self.vae.disable_slicing()
|
||||
|
||||
def enable_sequential_cpu_offload(self, gpu_id=0):
|
||||
r"""
|
||||
Offloads all models to CPU using accelerate, significantly reducing memory usage. When called, unet,
|
||||
@@ -191,10 +244,15 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
|
||||
device = torch.device(f"cuda:{gpu_id}")
|
||||
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae, self.safety_checker]:
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae]:
|
||||
if cpu_offloaded_model is not None:
|
||||
cpu_offload(cpu_offloaded_model, device)
|
||||
|
||||
if self.safety_checker is not None:
|
||||
# TODO(Patrick) - there is currently a bug with cpu offload of nn.Parameter in accelerate
|
||||
# fix by only offloading self.safety_checker for now
|
||||
cpu_offload(self.safety_checker.vision_model, device)
|
||||
|
||||
@property
|
||||
def _execution_device(self):
|
||||
r"""
|
||||
@@ -369,7 +427,7 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
)
|
||||
|
||||
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
|
||||
shape = (batch_size, num_channels_latents, height // 8, width // 8)
|
||||
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor)
|
||||
if latents is None:
|
||||
if device.type == "mps":
|
||||
# randn does not work reproducibly on mps
|
||||
@@ -389,8 +447,8 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -402,7 +460,6 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
return_dict: bool = True,
|
||||
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
|
||||
callback_steps: Optional[int] = 1,
|
||||
**kwargs,
|
||||
):
|
||||
r"""
|
||||
Function invoked when calling the pipeline for generation.
|
||||
@@ -410,9 +467,9 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -458,6 +515,9 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
# 1. Check inputs. Raise error if not correct
|
||||
self.check_inputs(prompt, height, width, callback_steps)
|
||||
@@ -496,25 +556,29 @@ class StableDiffusionPipeline(DiffusionPipeline):
|
||||
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
|
||||
|
||||
# 7. Denoising loop
|
||||
for i, t in enumerate(self.progress_bar(timesteps)):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
for i, t in enumerate(timesteps):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
|
||||
# call the callback, if provided
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
# call the callback, if provided
|
||||
if (i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0:
|
||||
progress_bar.update()
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
|
||||
# 8. Post-processing
|
||||
image = self.decode_latents(latents)
|
||||
|
||||
@@ -19,8 +19,10 @@ import torch
|
||||
|
||||
import PIL
|
||||
from diffusers.utils import is_accelerate_available
|
||||
from packaging import version
|
||||
from transformers import CLIPFeatureExtractor, CLIPVisionModelWithProjection
|
||||
|
||||
from ...configuration_utils import FrozenDict
|
||||
from ...models import AutoencoderKL, UNet2DConditionModel
|
||||
from ...pipeline_utils import DiffusionPipeline
|
||||
from ...schedulers import (
|
||||
@@ -31,7 +33,7 @@ from ...schedulers import (
|
||||
LMSDiscreteScheduler,
|
||||
PNDMScheduler,
|
||||
)
|
||||
from ...utils import logging
|
||||
from ...utils import deprecate, logging
|
||||
from . import StableDiffusionPipelineOutput
|
||||
from .safety_checker import StableDiffusionSafetyChecker
|
||||
|
||||
@@ -63,6 +65,7 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
feature_extractor ([`CLIPFeatureExtractor`]):
|
||||
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
|
||||
"""
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
@@ -79,10 +82,11 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
],
|
||||
safety_checker: StableDiffusionSafetyChecker,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warn(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
@@ -92,6 +96,33 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
|
||||
version.parse(unet.config._diffusers_version).base_version
|
||||
) < version.parse("0.9.0.dev0")
|
||||
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
|
||||
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
|
||||
deprecation_message = (
|
||||
"The configuration file of the unet has set the default `sample_size` to smaller than"
|
||||
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
|
||||
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
|
||||
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
|
||||
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
|
||||
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
|
||||
" in the config might lead to incorrect results in future versions. If you have downloaded this"
|
||||
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
|
||||
" the `unet/config.json` file"
|
||||
)
|
||||
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
|
||||
new_config = dict(unet.config)
|
||||
new_config["sample_size"] = 64
|
||||
unet._internal_dict = FrozenDict(new_config)
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
image_encoder=image_encoder,
|
||||
@@ -100,6 +131,8 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_xformers_memory_efficient_attention
|
||||
def enable_xformers_memory_efficient_attention(self):
|
||||
@@ -136,9 +169,14 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
if isinstance(self.unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.unet.config.attention_head_dim)
|
||||
|
||||
self.unet.set_attention_slice(slice_size)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_attention_slicing
|
||||
@@ -272,7 +310,7 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_latents
|
||||
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
|
||||
shape = (batch_size, num_channels_latents, height // 8, width // 8)
|
||||
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor)
|
||||
if latents is None:
|
||||
if device.type == "mps":
|
||||
# randn does not work reproducibly on mps
|
||||
@@ -292,8 +330,8 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
image: Union[PIL.Image.Image, List[PIL.Image.Image], torch.FloatTensor],
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
num_images_per_prompt: Optional[int] = 1,
|
||||
@@ -304,7 +342,6 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
return_dict: bool = True,
|
||||
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
|
||||
callback_steps: Optional[int] = 1,
|
||||
**kwargs,
|
||||
):
|
||||
r"""
|
||||
Function invoked when calling the pipeline for generation.
|
||||
@@ -315,9 +352,9 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
configuration of
|
||||
[this](https://huggingface.co/lambdalabs/sd-image-variations-diffusers/blob/main/feature_extractor/preprocessor_config.json)
|
||||
`CLIPFeatureExtractor`
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -360,6 +397,9 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
# 1. Check inputs. Raise error if not correct
|
||||
self.check_inputs(image, height, width, callback_steps)
|
||||
@@ -401,25 +441,29 @@ class StableDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
|
||||
|
||||
# 7. Denoising loop
|
||||
for i, t in enumerate(self.progress_bar(timesteps)):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
for i, t in enumerate(timesteps):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=image_embeddings).sample
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=image_embeddings).sample
|
||||
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
|
||||
# call the callback, if provided
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
# call the callback, if provided
|
||||
if (i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0:
|
||||
progress_bar.update()
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
|
||||
# 8. Post-processing
|
||||
image = self.decode_latents(latents)
|
||||
|
||||
@@ -20,6 +20,7 @@ import torch
|
||||
|
||||
import PIL
|
||||
from diffusers.utils import is_accelerate_available
|
||||
from packaging import version
|
||||
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
|
||||
|
||||
from ...configuration_utils import FrozenDict
|
||||
@@ -78,6 +79,7 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
feature_extractor ([`CLIPFeatureExtractor`]):
|
||||
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
|
||||
"""
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.__init__
|
||||
def __init__(
|
||||
@@ -96,6 +98,7 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
],
|
||||
safety_checker: StableDiffusionSafetyChecker,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -126,7 +129,7 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
new_config["clip_sample"] = False
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
@@ -136,6 +139,33 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
|
||||
version.parse(unet.config._diffusers_version).base_version
|
||||
) < version.parse("0.9.0.dev0")
|
||||
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
|
||||
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
|
||||
deprecation_message = (
|
||||
"The configuration file of the unet has set the default `sample_size` to smaller than"
|
||||
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
|
||||
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
|
||||
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
|
||||
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
|
||||
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
|
||||
" in the config might lead to incorrect results in future versions. If you have downloaded this"
|
||||
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
|
||||
" the `unet/config.json` file"
|
||||
)
|
||||
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
|
||||
new_config = dict(unet.config)
|
||||
new_config["sample_size"] = 64
|
||||
unet._internal_dict = FrozenDict(new_config)
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
text_encoder=text_encoder,
|
||||
@@ -145,6 +175,8 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_attention_slicing
|
||||
def enable_attention_slicing(self, slice_size: Optional[Union[str, int]] = "auto"):
|
||||
@@ -161,9 +193,14 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
if isinstance(self.unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.unet.config.attention_head_dim)
|
||||
|
||||
self.unet.set_attention_slice(slice_size)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_attention_slicing
|
||||
@@ -189,10 +226,15 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
|
||||
device = torch.device(f"cuda:{gpu_id}")
|
||||
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae, self.safety_checker]:
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae]:
|
||||
if cpu_offloaded_model is not None:
|
||||
cpu_offload(cpu_offloaded_model, device)
|
||||
|
||||
if self.safety_checker is not None:
|
||||
# TODO(Patrick) - there is currently a bug with cpu offload of nn.Parameter in accelerate
|
||||
# fix by only offloading self.safety_checker for now
|
||||
cpu_offload(self.safety_checker.vision_model, device)
|
||||
|
||||
@property
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline._execution_device
|
||||
def _execution_device(self):
|
||||
@@ -400,11 +442,11 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
t_start = max(num_inference_steps - init_timestep + offset, 0)
|
||||
timesteps = self.scheduler.timesteps[t_start:]
|
||||
|
||||
return timesteps
|
||||
return timesteps, num_inference_steps - t_start
|
||||
|
||||
def prepare_latents(self, init_image, timestep, batch_size, num_images_per_prompt, dtype, device, generator=None):
|
||||
init_image = init_image.to(device=device, dtype=dtype)
|
||||
init_latent_dist = self.vae.encode(init_image).latent_dist
|
||||
def prepare_latents(self, image, timestep, batch_size, num_images_per_prompt, dtype, device, generator=None):
|
||||
image = image.to(device=device, dtype=dtype)
|
||||
init_latent_dist = self.vae.encode(image).latent_dist
|
||||
init_latents = init_latent_dist.sample(generator=generator)
|
||||
init_latents = 0.18215 * init_latents
|
||||
|
||||
@@ -412,16 +454,16 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
# expand init_latents for batch_size
|
||||
deprecation_message = (
|
||||
f"You have passed {batch_size} text prompts (`prompt`), but only {init_latents.shape[0]} initial"
|
||||
" images (`init_image`). Initial images are now duplicating to match the number of text prompts. Note"
|
||||
" images (`image`). Initial images are now duplicating to match the number of text prompts. Note"
|
||||
" that this behavior is deprecated and will be removed in a version 1.0.0. Please make sure to update"
|
||||
" your script to pass as many init images as text prompts to suppress this warning."
|
||||
" your script to pass as many initial images as text prompts to suppress this warning."
|
||||
)
|
||||
deprecate("len(prompt) != len(init_image)", "1.0.0", deprecation_message, standard_warn=False)
|
||||
deprecate("len(prompt) != len(image)", "1.0.0", deprecation_message, standard_warn=False)
|
||||
additional_image_per_prompt = batch_size // init_latents.shape[0]
|
||||
init_latents = torch.cat([init_latents] * additional_image_per_prompt * num_images_per_prompt, dim=0)
|
||||
elif batch_size > init_latents.shape[0] and batch_size % init_latents.shape[0] != 0:
|
||||
raise ValueError(
|
||||
f"Cannot duplicate `init_image` of batch size {init_latents.shape[0]} to {batch_size} text prompts."
|
||||
f"Cannot duplicate `image` of batch size {init_latents.shape[0]} to {batch_size} text prompts."
|
||||
)
|
||||
else:
|
||||
init_latents = torch.cat([init_latents] * num_images_per_prompt, dim=0)
|
||||
@@ -439,7 +481,7 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
init_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
strength: float = 0.8,
|
||||
num_inference_steps: Optional[int] = 50,
|
||||
guidance_scale: Optional[float] = 7.5,
|
||||
@@ -459,15 +501,15 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
init_image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process.
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
|
||||
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
|
||||
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
|
||||
noise will be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
|
||||
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1. `image`
|
||||
will be used as a starting point, adding more noise to it the larger the `strength`. The number of
|
||||
denoising steps depends on the amount of noise initially added. When `strength` is 1, added noise will
|
||||
be maximum and the denoising process will run for the full number of iterations specified in
|
||||
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
expense of slower inference. This parameter will be modulated by `strength`.
|
||||
@@ -508,6 +550,10 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
message = "Please use `image` instead of `init_image`."
|
||||
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
|
||||
image = init_image or image
|
||||
|
||||
# 1. Check inputs
|
||||
self.check_inputs(prompt, strength, callback_steps)
|
||||
|
||||
@@ -525,42 +571,46 @@ class StableDiffusionImg2ImgPipeline(DiffusionPipeline):
|
||||
)
|
||||
|
||||
# 4. Preprocess image
|
||||
if isinstance(init_image, PIL.Image.Image):
|
||||
init_image = preprocess(init_image)
|
||||
if isinstance(image, PIL.Image.Image):
|
||||
image = preprocess(image)
|
||||
|
||||
# 5. set timesteps
|
||||
self.scheduler.set_timesteps(num_inference_steps, device=device)
|
||||
timesteps = self.get_timesteps(num_inference_steps, strength, device)
|
||||
timesteps, num_inference_steps = self.get_timesteps(num_inference_steps, strength, device)
|
||||
latent_timestep = timesteps[:1].repeat(batch_size * num_images_per_prompt)
|
||||
|
||||
# 6. Prepare latent variables
|
||||
latents = self.prepare_latents(
|
||||
init_image, latent_timestep, batch_size, num_images_per_prompt, text_embeddings.dtype, device, generator
|
||||
image, latent_timestep, batch_size, num_images_per_prompt, text_embeddings.dtype, device, generator
|
||||
)
|
||||
|
||||
# 7. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
|
||||
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
|
||||
|
||||
# 8. Denoising loop
|
||||
for i, t in enumerate(self.progress_bar(timesteps)):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
for i, t in enumerate(timesteps):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
|
||||
# call the callback, if provided
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
# call the callback, if provided
|
||||
if (i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0:
|
||||
progress_bar.update()
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
|
||||
# 9. Post-processing
|
||||
image = self.decode_latents(latents)
|
||||
|
||||
@@ -20,6 +20,7 @@ import torch
|
||||
|
||||
import PIL
|
||||
from diffusers.utils import is_accelerate_available
|
||||
from packaging import version
|
||||
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
|
||||
|
||||
from ...configuration_utils import FrozenDict
|
||||
@@ -150,6 +151,7 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
feature_extractor ([`CLIPFeatureExtractor`]):
|
||||
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
|
||||
"""
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
@@ -160,6 +162,7 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
|
||||
safety_checker: StableDiffusionSafetyChecker,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -191,7 +194,7 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
new_config["skip_prk_steps"] = True
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
@@ -201,6 +204,33 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
|
||||
version.parse(unet.config._diffusers_version).base_version
|
||||
) < version.parse("0.9.0.dev0")
|
||||
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
|
||||
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
|
||||
deprecation_message = (
|
||||
"The configuration file of the unet has set the default `sample_size` to smaller than"
|
||||
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
|
||||
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
|
||||
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
|
||||
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
|
||||
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
|
||||
" in the config might lead to incorrect results in future versions. If you have downloaded this"
|
||||
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
|
||||
" the `unet/config.json` file"
|
||||
)
|
||||
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
|
||||
new_config = dict(unet.config)
|
||||
new_config["sample_size"] = 64
|
||||
unet._internal_dict = FrozenDict(new_config)
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
text_encoder=text_encoder,
|
||||
@@ -210,6 +240,8 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_attention_slicing
|
||||
def enable_attention_slicing(self, slice_size: Optional[Union[str, int]] = "auto"):
|
||||
@@ -226,9 +258,14 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
if isinstance(self.unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.unet.config.attention_head_dim)
|
||||
|
||||
self.unet.set_attention_slice(slice_size)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_attention_slicing
|
||||
@@ -254,10 +291,15 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
|
||||
device = torch.device(f"cuda:{gpu_id}")
|
||||
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae, self.safety_checker]:
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae]:
|
||||
if cpu_offloaded_model is not None:
|
||||
cpu_offload(cpu_offloaded_model, device)
|
||||
|
||||
if self.safety_checker is not None:
|
||||
# TODO(Patrick) - there is currently a bug with cpu offload of nn.Parameter in accelerate
|
||||
# fix by only offloading self.safety_checker for now
|
||||
cpu_offload(self.safety_checker.vision_model, device)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_xformers_memory_efficient_attention
|
||||
def enable_xformers_memory_efficient_attention(self):
|
||||
r"""
|
||||
@@ -459,7 +501,7 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_latents
|
||||
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
|
||||
shape = (batch_size, num_channels_latents, height // 8, width // 8)
|
||||
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor)
|
||||
if latents is None:
|
||||
if device.type == "mps":
|
||||
# randn does not work reproducibly on mps
|
||||
@@ -481,7 +523,9 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
# resize the mask to latents shape as we concatenate the mask to the latents
|
||||
# we do that before converting to dtype to avoid breaking in case we're using cpu_offload
|
||||
# and half precision
|
||||
mask = torch.nn.functional.interpolate(mask, size=(height // 8, width // 8))
|
||||
mask = torch.nn.functional.interpolate(
|
||||
mask, size=(height // self.vae_scale_factor, width // self.vae_scale_factor)
|
||||
)
|
||||
mask = mask.to(device=device, dtype=dtype)
|
||||
|
||||
masked_image = masked_image.to(device=device, dtype=dtype)
|
||||
@@ -509,8 +553,8 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
prompt: Union[str, List[str]],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
mask_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -522,7 +566,6 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
return_dict: bool = True,
|
||||
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
|
||||
callback_steps: Optional[int] = 1,
|
||||
**kwargs,
|
||||
):
|
||||
r"""
|
||||
Function invoked when calling the pipeline for generation.
|
||||
@@ -538,9 +581,9 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
repainted, while black pixels will be preserved. If `mask_image` is a PIL image, it will be converted
|
||||
to a single channel (luminance) before use. If it's a tensor, it should contain one color channel (L)
|
||||
instead of 3, so the expected shape would be `(B, H, W, 1)`.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -586,6 +629,9 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
# 1. Check inputs
|
||||
self.check_inputs(prompt, height, width, callback_steps)
|
||||
@@ -609,7 +655,7 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
|
||||
# 5. set timesteps
|
||||
self.scheduler.set_timesteps(num_inference_steps, device=device)
|
||||
timesteps_tensor = self.scheduler.timesteps
|
||||
timesteps = self.scheduler.timesteps
|
||||
|
||||
# 6. Prepare latent variables
|
||||
num_channels_latents = self.vae.config.latent_channels
|
||||
@@ -653,28 +699,32 @@ class StableDiffusionInpaintPipeline(DiffusionPipeline):
|
||||
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
|
||||
|
||||
# 10. Denoising loop
|
||||
for i, t in enumerate(self.progress_bar(timesteps_tensor)):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
for i, t in enumerate(timesteps):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
|
||||
# concat latents, mask, masked_image_latents in the channel dimension
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
latent_model_input = torch.cat([latent_model_input, mask, masked_image_latents], dim=1)
|
||||
# concat latents, mask, masked_image_latents in the channel dimension
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
latent_model_input = torch.cat([latent_model_input, mask, masked_image_latents], dim=1)
|
||||
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
|
||||
# call the callback, if provided
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
# call the callback, if provided
|
||||
if (i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0:
|
||||
progress_bar.update()
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
|
||||
# 11. Post-processing
|
||||
image = self.decode_latents(latents)
|
||||
|
||||
@@ -20,6 +20,7 @@ import torch
|
||||
|
||||
import PIL
|
||||
from diffusers.utils import is_accelerate_available
|
||||
from packaging import version
|
||||
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
|
||||
|
||||
from ...configuration_utils import FrozenDict
|
||||
@@ -51,11 +52,11 @@ def preprocess_image(image):
|
||||
return 2.0 * image - 1.0
|
||||
|
||||
|
||||
def preprocess_mask(mask):
|
||||
def preprocess_mask(mask, scale_factor=8):
|
||||
mask = mask.convert("L")
|
||||
w, h = mask.size
|
||||
w, h = map(lambda x: x - x % 32, (w, h)) # resize to integer multiple of 32
|
||||
mask = mask.resize((w // 8, h // 8), resample=PIL_INTERPOLATION["nearest"])
|
||||
mask = mask.resize((w // scale_factor, h // scale_factor), resample=PIL_INTERPOLATION["nearest"])
|
||||
mask = np.array(mask).astype(np.float32) / 255.0
|
||||
mask = np.tile(mask, (4, 1, 1))
|
||||
mask = mask[None].transpose(0, 1, 2, 3) # what does this step do?
|
||||
@@ -91,6 +92,7 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
feature_extractor ([`CLIPFeatureExtractor`]):
|
||||
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
|
||||
"""
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.__init__
|
||||
def __init__(
|
||||
@@ -109,6 +111,7 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
],
|
||||
safety_checker: StableDiffusionSafetyChecker,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -139,7 +142,7 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
new_config["clip_sample"] = False
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
@@ -149,6 +152,33 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
|
||||
version.parse(unet.config._diffusers_version).base_version
|
||||
) < version.parse("0.9.0.dev0")
|
||||
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
|
||||
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
|
||||
deprecation_message = (
|
||||
"The configuration file of the unet has set the default `sample_size` to smaller than"
|
||||
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
|
||||
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
|
||||
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
|
||||
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
|
||||
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
|
||||
" in the config might lead to incorrect results in future versions. If you have downloaded this"
|
||||
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
|
||||
" the `unet/config.json` file"
|
||||
)
|
||||
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
|
||||
new_config = dict(unet.config)
|
||||
new_config["sample_size"] = 64
|
||||
unet._internal_dict = FrozenDict(new_config)
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
text_encoder=text_encoder,
|
||||
@@ -158,6 +188,8 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
safety_checker=safety_checker,
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_attention_slicing
|
||||
def enable_attention_slicing(self, slice_size: Optional[Union[str, int]] = "auto"):
|
||||
@@ -174,9 +206,14 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
if isinstance(self.unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.unet.config.attention_head_dim)
|
||||
|
||||
self.unet.set_attention_slice(slice_size)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_attention_slicing
|
||||
@@ -202,10 +239,15 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
|
||||
device = torch.device(f"cuda:{gpu_id}")
|
||||
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae, self.safety_checker]:
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae]:
|
||||
if cpu_offloaded_model is not None:
|
||||
cpu_offload(cpu_offloaded_model, device)
|
||||
|
||||
if self.safety_checker is not None:
|
||||
# TODO(Patrick) - there is currently a bug with cpu offload of nn.Parameter in accelerate
|
||||
# fix by only offloading self.safety_checker for now
|
||||
cpu_offload(self.safety_checker.vision_model, device)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_xformers_memory_efficient_attention
|
||||
def enable_xformers_memory_efficient_attention(self):
|
||||
r"""
|
||||
@@ -415,11 +457,11 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
t_start = max(num_inference_steps - init_timestep + offset, 0)
|
||||
timesteps = self.scheduler.timesteps[t_start:]
|
||||
|
||||
return timesteps
|
||||
return timesteps, num_inference_steps - t_start
|
||||
|
||||
def prepare_latents(self, init_image, timestep, batch_size, num_images_per_prompt, dtype, device, generator):
|
||||
init_image = init_image.to(device=self.device, dtype=dtype)
|
||||
init_latent_dist = self.vae.encode(init_image).latent_dist
|
||||
def prepare_latents(self, image, timestep, batch_size, num_images_per_prompt, dtype, device, generator):
|
||||
image = image.to(device=self.device, dtype=dtype)
|
||||
init_latent_dist = self.vae.encode(image).latent_dist
|
||||
init_latents = init_latent_dist.sample(generator=generator)
|
||||
init_latents = 0.18215 * init_latents
|
||||
|
||||
@@ -437,7 +479,7 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
init_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
mask_image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
strength: float = 0.8,
|
||||
num_inference_steps: Optional[int] = 50,
|
||||
@@ -458,19 +500,19 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
init_image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, that will be used as the starting point for the
|
||||
process. This is the image whose masked region will be inpainted.
|
||||
mask_image (`torch.FloatTensor` or `PIL.Image.Image`):
|
||||
`Image`, or tensor representing an image batch, to mask `init_image`. White pixels in the mask will be
|
||||
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be
|
||||
replaced by noise and therefore repainted, while black pixels will be preserved. If `mask_image` is a
|
||||
PIL image, it will be converted to a single channel (luminance) before use. If it's a tensor, it should
|
||||
contain one color channel (L) instead of 3, so the expected shape would be `(B, H, W, 1)`.
|
||||
strength (`float`, *optional*, defaults to 0.8):
|
||||
Conceptually, indicates how much to inpaint the masked area. Must be between 0 and 1. When `strength`
|
||||
is 1, the denoising process will be run on the masked area for the full number of iterations specified
|
||||
in `num_inference_steps`. `init_image` will be used as a reference for the masked area, adding more
|
||||
noise to that region the larger the `strength`. If `strength` is 0, no inpainting will occur.
|
||||
in `num_inference_steps`. `image` will be used as a reference for the masked area, adding more noise to
|
||||
that region the larger the `strength`. If `strength` is 0, no inpainting will occur.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The reference number of denoising steps. More denoising steps usually lead to a higher quality image at
|
||||
the expense of slower inference. This parameter will be modulated by `strength`, as explained above.
|
||||
@@ -511,6 +553,10 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
message = "Please use `image` instead of `init_image`."
|
||||
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
|
||||
image = init_image or image
|
||||
|
||||
# 1. Check inputs
|
||||
self.check_inputs(prompt, strength, callback_steps)
|
||||
|
||||
@@ -528,21 +574,21 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
)
|
||||
|
||||
# 4. Preprocess image and mask
|
||||
if not isinstance(init_image, torch.FloatTensor):
|
||||
init_image = preprocess_image(init_image)
|
||||
if not isinstance(image, torch.FloatTensor):
|
||||
image = preprocess_image(image)
|
||||
|
||||
if not isinstance(mask_image, torch.FloatTensor):
|
||||
mask_image = preprocess_mask(mask_image)
|
||||
mask_image = preprocess_mask(mask_image, self.vae_scale_factor)
|
||||
|
||||
# 5. set timesteps
|
||||
self.scheduler.set_timesteps(num_inference_steps, device=device)
|
||||
timesteps = self.get_timesteps(num_inference_steps, strength, device)
|
||||
timesteps, num_inference_steps = self.get_timesteps(num_inference_steps, strength, device)
|
||||
latent_timestep = timesteps[:1].repeat(batch_size * num_images_per_prompt)
|
||||
|
||||
# 6. Prepare latent variables
|
||||
# encode the init image into latents and scale the latents
|
||||
latents, init_latents_orig, noise = self.prepare_latents(
|
||||
init_image, latent_timestep, batch_size, num_images_per_prompt, text_embeddings.dtype, device, generator
|
||||
image, latent_timestep, batch_size, num_images_per_prompt, text_embeddings.dtype, device, generator
|
||||
)
|
||||
|
||||
# 7. Prepare mask latent
|
||||
@@ -553,29 +599,33 @@ class StableDiffusionInpaintPipelineLegacy(DiffusionPipeline):
|
||||
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
|
||||
|
||||
# 9. Denoising loop
|
||||
for i, t in enumerate(self.progress_bar(timesteps)):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
for i, t in enumerate(timesteps):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
# masking
|
||||
init_latents_proper = self.scheduler.add_noise(init_latents_orig, noise, torch.tensor([t]))
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
# masking
|
||||
init_latents_proper = self.scheduler.add_noise(init_latents_orig, noise, torch.tensor([t]))
|
||||
|
||||
latents = (init_latents_proper * mask) + (latents * (1 - mask))
|
||||
latents = (init_latents_proper * mask) + (latents * (1 - mask))
|
||||
|
||||
# call the callback, if provided
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
# call the callback, if provided
|
||||
if (i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0:
|
||||
progress_bar.update()
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
|
||||
# 10. Post-processing
|
||||
image = self.decode_latents(latents)
|
||||
|
||||
@@ -0,0 +1,555 @@
|
||||
# Copyright 2022 The HuggingFace Team. All rights reserved.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License.
|
||||
|
||||
import inspect
|
||||
from typing import Callable, List, Optional, Union
|
||||
|
||||
import numpy as np
|
||||
import torch
|
||||
|
||||
import PIL
|
||||
from diffusers.utils import is_accelerate_available
|
||||
from transformers import CLIPTextModel, CLIPTokenizer
|
||||
|
||||
from ...models import AutoencoderKL, UNet2DConditionModel
|
||||
from ...pipeline_utils import DiffusionPipeline, ImagePipelineOutput
|
||||
from ...schedulers import DDIMScheduler, DDPMScheduler, LMSDiscreteScheduler, PNDMScheduler
|
||||
from ...utils import logging
|
||||
|
||||
|
||||
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
|
||||
|
||||
|
||||
def preprocess(image):
|
||||
# resize to multiple of 64
|
||||
width, height = image.size
|
||||
width = width - width % 64
|
||||
height = height - height % 64
|
||||
image = image.resize((width, height))
|
||||
|
||||
image = np.array(image.convert("RGB"))
|
||||
image = image[None].transpose(0, 3, 1, 2)
|
||||
image = torch.from_numpy(image).to(dtype=torch.float32) / 127.5 - 1.0
|
||||
return image
|
||||
|
||||
|
||||
class StableDiffusionUpscalePipeline(DiffusionPipeline):
|
||||
r"""
|
||||
Pipeline for text-guided image super-resolution using Stable Diffusion 2.
|
||||
|
||||
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods the
|
||||
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
|
||||
|
||||
Args:
|
||||
vae ([`AutoencoderKL`]):
|
||||
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
|
||||
text_encoder ([`CLIPTextModel`]):
|
||||
Frozen text-encoder. Stable Diffusion uses the text portion of
|
||||
[CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModel), specifically
|
||||
the [clip-vit-large-patch14](https://huggingface.co/openai/clip-vit-large-patch14) variant.
|
||||
tokenizer (`CLIPTokenizer`):
|
||||
Tokenizer of class
|
||||
[CLIPTokenizer](https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer).
|
||||
unet ([`UNet2DConditionModel`]): Conditional U-Net architecture to denoise the encoded image latents.
|
||||
low_res_scheduler ([`SchedulerMixin`]):
|
||||
A scheduler used to add initial noise to the low res conditioning image. It must be an instance of
|
||||
[`DDPMScheduler`].
|
||||
scheduler ([`SchedulerMixin`]):
|
||||
A scheduler to be used in combination with `unet` to denoise the encoded image latents. Can be one of
|
||||
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
|
||||
"""
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
vae: AutoencoderKL,
|
||||
text_encoder: CLIPTextModel,
|
||||
tokenizer: CLIPTokenizer,
|
||||
unet: UNet2DConditionModel,
|
||||
low_res_scheduler: DDPMScheduler,
|
||||
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
|
||||
max_noise_level: int = 350,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
text_encoder=text_encoder,
|
||||
tokenizer=tokenizer,
|
||||
unet=unet,
|
||||
low_res_scheduler=low_res_scheduler,
|
||||
scheduler=scheduler,
|
||||
)
|
||||
self.register_to_config(max_noise_level=max_noise_level)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_attention_slicing
|
||||
def enable_attention_slicing(self, slice_size: Optional[Union[str, int]] = "auto"):
|
||||
r"""
|
||||
Enable sliced attention computation.
|
||||
|
||||
When this option is enabled, the attention module will split the input tensor in slices, to compute attention
|
||||
in several steps. This is useful to save some memory in exchange for a small speed decrease.
|
||||
|
||||
Args:
|
||||
slice_size (`str` or `int`, *optional*, defaults to `"auto"`):
|
||||
When `"auto"`, halves the input to the attention heads, so attention will be computed in two steps. If
|
||||
a number is provided, uses as many slices as `attention_head_dim // slice_size`. In this case,
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
if isinstance(self.unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.unet.config.attention_head_dim)
|
||||
|
||||
self.unet.set_attention_slice(slice_size)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_attention_slicing
|
||||
def disable_attention_slicing(self):
|
||||
r"""
|
||||
Disable sliced attention computation. If `enable_attention_slicing` was previously invoked, this method will go
|
||||
back to computing attention in one step.
|
||||
"""
|
||||
# set slice_size = `None` to disable `attention slicing`
|
||||
self.enable_attention_slicing(None)
|
||||
|
||||
def enable_sequential_cpu_offload(self, gpu_id=0):
|
||||
r"""
|
||||
Offloads all models to CPU using accelerate, significantly reducing memory usage. When called, unet,
|
||||
text_encoder, vae and safety checker have their state dicts saved to CPU and then are moved to a
|
||||
`torch.device('meta') and loaded to GPU only when their specific submodule has its `forward` method called.
|
||||
"""
|
||||
if is_accelerate_available():
|
||||
from accelerate import cpu_offload
|
||||
else:
|
||||
raise ImportError("Please install accelerate via `pip install accelerate`")
|
||||
|
||||
device = torch.device(f"cuda:{gpu_id}")
|
||||
|
||||
for cpu_offloaded_model in [self.unet, self.text_encoder]:
|
||||
if cpu_offloaded_model is not None:
|
||||
cpu_offload(cpu_offloaded_model, device)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_xformers_memory_efficient_attention
|
||||
def enable_xformers_memory_efficient_attention(self):
|
||||
r"""
|
||||
Enable memory efficient attention as implemented in xformers.
|
||||
|
||||
When this option is enabled, you should observe lower GPU memory usage and a potential speed up at inference
|
||||
time. Speed up at training time is not guaranteed.
|
||||
|
||||
Warning: When Memory Efficient Attention and Sliced attention are both enabled, the Memory Efficient Attention
|
||||
is used.
|
||||
"""
|
||||
self.unet.set_use_memory_efficient_attention_xformers(True)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_xformers_memory_efficient_attention
|
||||
def disable_xformers_memory_efficient_attention(self):
|
||||
r"""
|
||||
Disable memory efficient attention as implemented in xformers.
|
||||
"""
|
||||
self.unet.set_use_memory_efficient_attention_xformers(False)
|
||||
|
||||
@property
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline._execution_device
|
||||
def _execution_device(self):
|
||||
r"""
|
||||
Returns the device on which the pipeline's models will be executed. After calling
|
||||
`pipeline.enable_sequential_cpu_offload()` the execution device can only be inferred from Accelerate's module
|
||||
hooks.
|
||||
"""
|
||||
if self.device != torch.device("meta") or not hasattr(self.unet, "_hf_hook"):
|
||||
return self.device
|
||||
for module in self.unet.modules():
|
||||
if (
|
||||
hasattr(module, "_hf_hook")
|
||||
and hasattr(module._hf_hook, "execution_device")
|
||||
and module._hf_hook.execution_device is not None
|
||||
):
|
||||
return torch.device(module._hf_hook.execution_device)
|
||||
return self.device
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline._encode_prompt
|
||||
def _encode_prompt(self, prompt, device, num_images_per_prompt, do_classifier_free_guidance, negative_prompt):
|
||||
r"""
|
||||
Encodes the prompt into text encoder hidden states.
|
||||
|
||||
Args:
|
||||
prompt (`str` or `list(int)`):
|
||||
prompt to be encoded
|
||||
device: (`torch.device`):
|
||||
torch device
|
||||
num_images_per_prompt (`int`):
|
||||
number of images that should be generated per prompt
|
||||
do_classifier_free_guidance (`bool`):
|
||||
whether to use classifier free guidance or not
|
||||
negative_prompt (`str` or `List[str]`):
|
||||
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
|
||||
if `guidance_scale` is less than `1`).
|
||||
"""
|
||||
batch_size = len(prompt) if isinstance(prompt, list) else 1
|
||||
|
||||
text_inputs = self.tokenizer(
|
||||
prompt,
|
||||
padding="max_length",
|
||||
max_length=self.tokenizer.model_max_length,
|
||||
truncation=True,
|
||||
return_tensors="pt",
|
||||
)
|
||||
text_input_ids = text_inputs.input_ids
|
||||
untruncated_ids = self.tokenizer(prompt, padding="max_length", return_tensors="pt").input_ids
|
||||
|
||||
if not torch.equal(text_input_ids, untruncated_ids):
|
||||
removed_text = self.tokenizer.batch_decode(untruncated_ids[:, self.tokenizer.model_max_length - 1 : -1])
|
||||
logger.warning(
|
||||
"The following part of your input was truncated because CLIP can only handle sequences up to"
|
||||
f" {self.tokenizer.model_max_length} tokens: {removed_text}"
|
||||
)
|
||||
|
||||
if hasattr(self.text_encoder.config, "use_attention_mask") and self.text_encoder.config.use_attention_mask:
|
||||
attention_mask = text_inputs.attention_mask.to(device)
|
||||
else:
|
||||
attention_mask = None
|
||||
|
||||
text_embeddings = self.text_encoder(
|
||||
text_input_ids.to(device),
|
||||
attention_mask=attention_mask,
|
||||
)
|
||||
text_embeddings = text_embeddings[0]
|
||||
|
||||
# duplicate text embeddings for each generation per prompt, using mps friendly method
|
||||
bs_embed, seq_len, _ = text_embeddings.shape
|
||||
text_embeddings = text_embeddings.repeat(1, num_images_per_prompt, 1)
|
||||
text_embeddings = text_embeddings.view(bs_embed * num_images_per_prompt, seq_len, -1)
|
||||
|
||||
# get unconditional embeddings for classifier free guidance
|
||||
if do_classifier_free_guidance:
|
||||
uncond_tokens: List[str]
|
||||
if negative_prompt is None:
|
||||
uncond_tokens = [""] * batch_size
|
||||
elif type(prompt) is not type(negative_prompt):
|
||||
raise TypeError(
|
||||
f"`negative_prompt` should be the same type to `prompt`, but got {type(negative_prompt)} !="
|
||||
f" {type(prompt)}."
|
||||
)
|
||||
elif isinstance(negative_prompt, str):
|
||||
uncond_tokens = [negative_prompt]
|
||||
elif batch_size != len(negative_prompt):
|
||||
raise ValueError(
|
||||
f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:"
|
||||
f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches"
|
||||
" the batch size of `prompt`."
|
||||
)
|
||||
else:
|
||||
uncond_tokens = negative_prompt
|
||||
|
||||
max_length = text_input_ids.shape[-1]
|
||||
uncond_input = self.tokenizer(
|
||||
uncond_tokens,
|
||||
padding="max_length",
|
||||
max_length=max_length,
|
||||
truncation=True,
|
||||
return_tensors="pt",
|
||||
)
|
||||
|
||||
if hasattr(self.text_encoder.config, "use_attention_mask") and self.text_encoder.config.use_attention_mask:
|
||||
attention_mask = uncond_input.attention_mask.to(device)
|
||||
else:
|
||||
attention_mask = None
|
||||
|
||||
uncond_embeddings = self.text_encoder(
|
||||
uncond_input.input_ids.to(device),
|
||||
attention_mask=attention_mask,
|
||||
)
|
||||
uncond_embeddings = uncond_embeddings[0]
|
||||
|
||||
# duplicate unconditional embeddings for each generation per prompt, using mps friendly method
|
||||
seq_len = uncond_embeddings.shape[1]
|
||||
uncond_embeddings = uncond_embeddings.repeat(1, num_images_per_prompt, 1)
|
||||
uncond_embeddings = uncond_embeddings.view(batch_size * num_images_per_prompt, seq_len, -1)
|
||||
|
||||
# For classifier free guidance, we need to do two forward passes.
|
||||
# Here we concatenate the unconditional and text embeddings into a single batch
|
||||
# to avoid doing two forward passes
|
||||
text_embeddings = torch.cat([uncond_embeddings, text_embeddings])
|
||||
|
||||
return text_embeddings
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_extra_step_kwargs
|
||||
def prepare_extra_step_kwargs(self, generator, eta):
|
||||
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature
|
||||
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers.
|
||||
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502
|
||||
# and should be between [0, 1]
|
||||
|
||||
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys())
|
||||
extra_step_kwargs = {}
|
||||
if accepts_eta:
|
||||
extra_step_kwargs["eta"] = eta
|
||||
|
||||
# check if the scheduler accepts generator
|
||||
accepts_generator = "generator" in set(inspect.signature(self.scheduler.step).parameters.keys())
|
||||
if accepts_generator:
|
||||
extra_step_kwargs["generator"] = generator
|
||||
return extra_step_kwargs
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.decode_latents with 0.18215->0.08333
|
||||
def decode_latents(self, latents):
|
||||
latents = 1 / 0.08333 * latents
|
||||
image = self.vae.decode(latents).sample
|
||||
image = (image / 2 + 0.5).clamp(0, 1)
|
||||
# we always cast to float32 as this does not cause significant overhead and is compatible with bfloa16
|
||||
image = image.cpu().permute(0, 2, 3, 1).float().numpy()
|
||||
return image
|
||||
|
||||
def check_inputs(self, prompt, image, noise_level, callback_steps):
|
||||
if not isinstance(prompt, str) and not isinstance(prompt, list):
|
||||
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}")
|
||||
|
||||
if (
|
||||
not isinstance(image, torch.Tensor)
|
||||
and not isinstance(image, PIL.Image.Image)
|
||||
and not isinstance(image, list)
|
||||
):
|
||||
raise ValueError(
|
||||
f"`image` has to be of type `torch.Tensor`, `PIL.Image.Image` or `list` but is {type(image)}"
|
||||
)
|
||||
|
||||
# verify batch size of prompt and image are same if image is a list or tensor
|
||||
if isinstance(image, list) or isinstance(image, torch.Tensor):
|
||||
if isinstance(prompt, str):
|
||||
batch_size = 1
|
||||
else:
|
||||
batch_size = len(prompt)
|
||||
if isinstance(image, list):
|
||||
image_batch_size = len(image)
|
||||
else:
|
||||
image_batch_size = image.shape[0]
|
||||
if batch_size != image_batch_size:
|
||||
raise ValueError(
|
||||
f"`prompt` has batch size {batch_size} and `image` has batch size {image_batch_size}."
|
||||
" Please make sure that passed `prompt` matches the batch size of `image`."
|
||||
)
|
||||
|
||||
# check noise level
|
||||
if noise_level > self.config.max_noise_level:
|
||||
raise ValueError(f"`noise_level` has to be <= {self.config.max_noise_level} but is {noise_level}")
|
||||
|
||||
if (callback_steps is None) or (
|
||||
callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0)
|
||||
):
|
||||
raise ValueError(
|
||||
f"`callback_steps` has to be a positive integer but is {callback_steps} of type"
|
||||
f" {type(callback_steps)}."
|
||||
)
|
||||
|
||||
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
|
||||
shape = (batch_size, num_channels_latents, height, width)
|
||||
if latents is None:
|
||||
if device.type == "mps":
|
||||
# randn does not work reproducibly on mps
|
||||
latents = torch.randn(shape, generator=generator, device="cpu", dtype=dtype).to(device)
|
||||
else:
|
||||
latents = torch.randn(shape, generator=generator, device=device, dtype=dtype)
|
||||
else:
|
||||
if latents.shape != shape:
|
||||
raise ValueError(f"Unexpected latents shape, got {latents.shape}, expected {shape}")
|
||||
latents = latents.to(device)
|
||||
|
||||
# scale the initial noise by the standard deviation required by the scheduler
|
||||
latents = latents * self.scheduler.init_noise_sigma
|
||||
return latents
|
||||
|
||||
@torch.no_grad()
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image, List[PIL.Image.Image]],
|
||||
num_inference_steps: int = 75,
|
||||
guidance_scale: float = 9.0,
|
||||
noise_level: int = 20,
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
num_images_per_prompt: Optional[int] = 1,
|
||||
eta: float = 0.0,
|
||||
generator: Optional[torch.Generator] = None,
|
||||
latents: Optional[torch.FloatTensor] = None,
|
||||
output_type: Optional[str] = "pil",
|
||||
return_dict: bool = True,
|
||||
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
|
||||
callback_steps: Optional[int] = 1,
|
||||
):
|
||||
r"""
|
||||
Function invoked when calling the pipeline for generation.
|
||||
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
image (`PIL.Image.Image` or List[`PIL.Image.Image`] or `torch.FloatTensor`):
|
||||
`Image`, or tensor representing an image batch which will be upscaled. *
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
expense of slower inference.
|
||||
guidance_scale (`float`, *optional*, defaults to 7.5):
|
||||
Guidance scale as defined in [Classifier-Free Diffusion Guidance](https://arxiv.org/abs/2207.12598).
|
||||
`guidance_scale` is defined as `w` of equation 2. of [Imagen
|
||||
Paper](https://arxiv.org/pdf/2205.11487.pdf). Guidance scale is enabled by setting `guidance_scale >
|
||||
1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`,
|
||||
usually at the expense of lower image quality.
|
||||
negative_prompt (`str` or `List[str]`, *optional*):
|
||||
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
|
||||
if `guidance_scale` is less than `1`).
|
||||
num_images_per_prompt (`int`, *optional*, defaults to 1):
|
||||
The number of images to generate per prompt.
|
||||
eta (`float`, *optional*, defaults to 0.0):
|
||||
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
|
||||
[`schedulers.DDIMScheduler`], will be ignored for others.
|
||||
generator (`torch.Generator`, *optional*):
|
||||
A [torch generator](https://pytorch.org/docs/stable/generated/torch.Generator.html) to make generation
|
||||
deterministic.
|
||||
latents (`torch.FloatTensor`, *optional*):
|
||||
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
|
||||
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
|
||||
tensor will ge generated by sampling using the supplied random `generator`.
|
||||
output_type (`str`, *optional*, defaults to `"pil"`):
|
||||
The output format of the generate image. Choose between
|
||||
[PIL](https://pillow.readthedocs.io/en/stable/): `PIL.Image.Image` or `np.array`.
|
||||
return_dict (`bool`, *optional*, defaults to `True`):
|
||||
Whether or not to return a [`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] instead of a
|
||||
plain tuple.
|
||||
callback (`Callable`, *optional*):
|
||||
A function that will be called every `callback_steps` steps during inference. The function will be
|
||||
called with the following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`.
|
||||
callback_steps (`int`, *optional*, defaults to 1):
|
||||
The frequency at which the `callback` function will be called. If not specified, the callback will be
|
||||
called at every step.
|
||||
|
||||
Returns:
|
||||
[`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] or `tuple`:
|
||||
[`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] if `return_dict` is True, otherwise a `tuple.
|
||||
When returning a tuple, the first element is a list with the generated images, and the second element is a
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
|
||||
# 1. Check inputs
|
||||
self.check_inputs(prompt, image, noise_level, callback_steps)
|
||||
|
||||
# 2. Define call parameters
|
||||
batch_size = 1 if isinstance(prompt, str) else len(prompt)
|
||||
device = self._execution_device
|
||||
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
|
||||
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
|
||||
# corresponds to doing no classifier free guidance.
|
||||
do_classifier_free_guidance = guidance_scale > 1.0
|
||||
|
||||
# 3. Encode input prompt
|
||||
text_embeddings = self._encode_prompt(
|
||||
prompt, device, num_images_per_prompt, do_classifier_free_guidance, negative_prompt
|
||||
)
|
||||
|
||||
# 4. Preprocess image
|
||||
image = [image] if isinstance(image, PIL.Image.Image) else image
|
||||
if isinstance(image, list):
|
||||
image = [preprocess(img) for img in image]
|
||||
image = torch.cat(image, dim=0)
|
||||
image = image.to(dtype=text_embeddings.dtype, device=device)
|
||||
|
||||
# 5. set timesteps
|
||||
self.scheduler.set_timesteps(num_inference_steps, device=device)
|
||||
timesteps = self.scheduler.timesteps
|
||||
|
||||
# 5. Add noise to image
|
||||
noise_level = torch.tensor([noise_level], dtype=torch.long, device=device)
|
||||
if device.type == "mps":
|
||||
# randn does not work reproducibly on mps
|
||||
noise = torch.randn(image.shape, generator=generator, device="cpu", dtype=text_embeddings.dtype).to(device)
|
||||
else:
|
||||
noise = torch.randn(image.shape, generator=generator, device=device, dtype=text_embeddings.dtype)
|
||||
image = self.low_res_scheduler.add_noise(image, noise, noise_level)
|
||||
image = torch.cat([image] * 2) if do_classifier_free_guidance else image
|
||||
noise_level = torch.cat([noise_level] * 2) if do_classifier_free_guidance else noise_level
|
||||
|
||||
# 6. Prepare latent variables
|
||||
height, width = image.shape[2:]
|
||||
num_channels_latents = self.vae.config.latent_channels
|
||||
latents = self.prepare_latents(
|
||||
batch_size * num_images_per_prompt,
|
||||
num_channels_latents,
|
||||
height,
|
||||
width,
|
||||
text_embeddings.dtype,
|
||||
device,
|
||||
generator,
|
||||
latents,
|
||||
)
|
||||
|
||||
# 7. Check that sizes of image and latents match
|
||||
num_channels_image = image.shape[1]
|
||||
if num_channels_latents + num_channels_image != self.unet.config.in_channels:
|
||||
raise ValueError(
|
||||
f"Incorrect configuration settings! The config of `pipeline.unet`: {self.unet.config} expects"
|
||||
f" {self.unet.config.in_channels} but received `num_channels_latents`: {num_channels_latents} +"
|
||||
f" `num_channels_image`: {num_channels_image} "
|
||||
f" = {num_channels_latents+num_channels_image}. Please verify the config of"
|
||||
" `pipeline.unet` or your `image` input."
|
||||
)
|
||||
|
||||
# 8. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
|
||||
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
|
||||
|
||||
# 9. Denoising loop
|
||||
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
for i, t in enumerate(timesteps):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
|
||||
|
||||
# concat latents, mask, masked_image_latents in the channel dimension
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
latent_model_input = torch.cat([latent_model_input, image], dim=1)
|
||||
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(
|
||||
latent_model_input, t, encoder_hidden_states=text_embeddings, class_labels=noise_level
|
||||
).sample
|
||||
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
|
||||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||||
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
|
||||
# call the callback, if provided
|
||||
if (i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0:
|
||||
progress_bar.update()
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
|
||||
# 10. Post-processing
|
||||
# make sure the VAE is in float32 mode, as it overflows in float16
|
||||
self.vae.to(dtype=torch.float32)
|
||||
image = self.decode_latents(latents.float())
|
||||
|
||||
# 11. Convert to PIL
|
||||
if output_type == "pil":
|
||||
image = self.numpy_to_pil(image)
|
||||
|
||||
if not return_dict:
|
||||
return (image,)
|
||||
|
||||
return ImagePipelineOutput(images=image)
|
||||
@@ -5,6 +5,7 @@ from typing import Callable, List, Optional, Union
|
||||
import numpy as np
|
||||
import torch
|
||||
|
||||
from packaging import version
|
||||
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
|
||||
|
||||
from ...configuration_utils import FrozenDict
|
||||
@@ -56,6 +57,8 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
|
||||
"""
|
||||
|
||||
_optional_components = ["safety_checker", "feature_extractor"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
vae: AutoencoderKL,
|
||||
@@ -72,6 +75,7 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
],
|
||||
safety_checker: SafeStableDiffusionSafetyChecker,
|
||||
feature_extractor: CLIPFeatureExtractor,
|
||||
requires_safety_checker: bool = True,
|
||||
):
|
||||
super().__init__()
|
||||
safety_concept: Optional[str] = (
|
||||
@@ -107,7 +111,7 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
new_config["clip_sample"] = False
|
||||
scheduler._internal_dict = FrozenDict(new_config)
|
||||
|
||||
if safety_checker is None:
|
||||
if safety_checker is None and requires_safety_checker:
|
||||
logger.warning(
|
||||
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
|
||||
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
|
||||
@@ -117,6 +121,33 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
|
||||
)
|
||||
|
||||
if safety_checker is not None and feature_extractor is None:
|
||||
raise ValueError(
|
||||
"Make sure to define a feature extractor when loading {self.__class__} if you want to use the safety"
|
||||
" checker. If you do not want to use the safety checker, you can pass `'safety_checker=None'` instead."
|
||||
)
|
||||
|
||||
is_unet_version_less_0_9_0 = hasattr(unet.config, "_diffusers_version") and version.parse(
|
||||
version.parse(unet.config._diffusers_version).base_version
|
||||
) < version.parse("0.9.0.dev0")
|
||||
is_unet_sample_size_less_64 = hasattr(unet.config, "sample_size") and unet.config.sample_size < 64
|
||||
if is_unet_version_less_0_9_0 and is_unet_sample_size_less_64:
|
||||
deprecation_message = (
|
||||
"The configuration file of the unet has set the default `sample_size` to smaller than"
|
||||
" 64 which seems highly unlikely .If you're checkpoint is a fine-tuned version of any of the"
|
||||
" following: \n- CompVis/stable-diffusion-v1-4 \n- CompVis/stable-diffusion-v1-3 \n-"
|
||||
" CompVis/stable-diffusion-v1-2 \n- CompVis/stable-diffusion-v1-1 \n- runwayml/stable-diffusion-v1-5"
|
||||
" \n- runwayml/stable-diffusion-inpainting \n you should change 'sample_size' to 64 in the"
|
||||
" configuration file. Please make sure to update the config accordingly as leaving `sample_size=32`"
|
||||
" in the config might lead to incorrect results in future versions. If you have downloaded this"
|
||||
" checkpoint from the Hugging Face Hub, it would be very nice if you could open a Pull request for"
|
||||
" the `unet/config.json` file"
|
||||
)
|
||||
deprecate("sample_size<64", "1.0.0", deprecation_message, standard_warn=False)
|
||||
new_config = dict(unet.config)
|
||||
new_config["sample_size"] = 64
|
||||
unet._internal_dict = FrozenDict(new_config)
|
||||
|
||||
self.register_modules(
|
||||
vae=vae,
|
||||
text_encoder=text_encoder,
|
||||
@@ -127,6 +158,8 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
feature_extractor=feature_extractor,
|
||||
)
|
||||
self._safety_text_concept = safety_concept
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
self.register_to_config(requires_safety_checker=requires_safety_checker)
|
||||
|
||||
@property
|
||||
def safety_concept(self):
|
||||
@@ -433,7 +466,7 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_latents
|
||||
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
|
||||
shape = (batch_size, num_channels_latents, height // 8, width // 8)
|
||||
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor)
|
||||
if latents is None:
|
||||
if device.type == "mps":
|
||||
# randn does not work reproducibly on mps
|
||||
@@ -495,8 +528,8 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -513,7 +546,6 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
sld_threshold: Optional[float] = 0.01,
|
||||
sld_momentum_scale: Optional[float] = 0.3,
|
||||
sld_mom_beta: Optional[float] = 0.4,
|
||||
**kwargs,
|
||||
):
|
||||
r"""
|
||||
Function invoked when calling the pipeline for generation.
|
||||
@@ -521,9 +553,9 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -589,6 +621,9 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
# 1. Check inputs. Raise error if not correct
|
||||
self.check_inputs(prompt, height, width, callback_steps)
|
||||
@@ -633,63 +668,71 @@ class StableDiffusionPipelineSafe(DiffusionPipeline):
|
||||
|
||||
safety_momentum = None
|
||||
|
||||
for i, t in enumerate(self.progress_bar(timesteps)):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = (
|
||||
torch.cat([latents] * (3 if enable_safety_guidance else 2)) if do_classifier_free_guidance else latents
|
||||
)
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
num_warmup_steps = len(timesteps) - num_inference_steps * self.scheduler.order
|
||||
with self.progress_bar(total=num_inference_steps) as progress_bar:
|
||||
for i, t in enumerate(timesteps):
|
||||
# expand the latents if we are doing classifier free guidance
|
||||
latent_model_input = (
|
||||
torch.cat([latents] * (3 if enable_safety_guidance else 2))
|
||||
if do_classifier_free_guidance
|
||||
else latents
|
||||
)
|
||||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||||
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
# predict the noise residual
|
||||
noise_pred = self.unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample
|
||||
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_out = noise_pred.chunk((3 if enable_safety_guidance else 2))
|
||||
noise_pred_uncond, noise_pred_text = noise_pred_out[0], noise_pred_out[1]
|
||||
# perform guidance
|
||||
if do_classifier_free_guidance:
|
||||
noise_pred_out = noise_pred.chunk((3 if enable_safety_guidance else 2))
|
||||
noise_pred_uncond, noise_pred_text = noise_pred_out[0], noise_pred_out[1]
|
||||
|
||||
# default classifier free guidance
|
||||
noise_guidance = noise_pred_text - noise_pred_uncond
|
||||
# default classifier free guidance
|
||||
noise_guidance = noise_pred_text - noise_pred_uncond
|
||||
|
||||
# Perform SLD guidance
|
||||
if enable_safety_guidance:
|
||||
if safety_momentum is None:
|
||||
safety_momentum = torch.zeros_like(noise_guidance)
|
||||
noise_pred_safety_concept = noise_pred_out[2]
|
||||
# Perform SLD guidance
|
||||
if enable_safety_guidance:
|
||||
if safety_momentum is None:
|
||||
safety_momentum = torch.zeros_like(noise_guidance)
|
||||
noise_pred_safety_concept = noise_pred_out[2]
|
||||
|
||||
# Equation 6
|
||||
scale = torch.clamp(
|
||||
torch.abs((noise_pred_text - noise_pred_safety_concept)) * sld_guidance_scale, max=1.0
|
||||
)
|
||||
# Equation 6
|
||||
scale = torch.clamp(
|
||||
torch.abs((noise_pred_text - noise_pred_safety_concept)) * sld_guidance_scale, max=1.0
|
||||
)
|
||||
|
||||
# Equation 6
|
||||
safety_concept_scale = torch.where(
|
||||
(noise_pred_text - noise_pred_safety_concept) >= sld_threshold, torch.zeros_like(scale), scale
|
||||
)
|
||||
# Equation 6
|
||||
safety_concept_scale = torch.where(
|
||||
(noise_pred_text - noise_pred_safety_concept) >= sld_threshold,
|
||||
torch.zeros_like(scale),
|
||||
scale,
|
||||
)
|
||||
|
||||
# Equation 4
|
||||
noise_guidance_safety = torch.mul(
|
||||
(noise_pred_safety_concept - noise_pred_uncond), safety_concept_scale
|
||||
)
|
||||
# Equation 4
|
||||
noise_guidance_safety = torch.mul(
|
||||
(noise_pred_safety_concept - noise_pred_uncond), safety_concept_scale
|
||||
)
|
||||
|
||||
# Equation 7
|
||||
noise_guidance_safety = noise_guidance_safety + sld_momentum_scale * safety_momentum
|
||||
# Equation 7
|
||||
noise_guidance_safety = noise_guidance_safety + sld_momentum_scale * safety_momentum
|
||||
|
||||
# Equation 8
|
||||
safety_momentum = sld_mom_beta * safety_momentum + (1 - sld_mom_beta) * noise_guidance_safety
|
||||
# Equation 8
|
||||
safety_momentum = sld_mom_beta * safety_momentum + (1 - sld_mom_beta) * noise_guidance_safety
|
||||
|
||||
if i >= sld_warmup_steps: # Warmup
|
||||
# Equation 3
|
||||
noise_guidance = noise_guidance - noise_guidance_safety
|
||||
if i >= sld_warmup_steps: # Warmup
|
||||
# Equation 3
|
||||
noise_guidance = noise_guidance - noise_guidance_safety
|
||||
|
||||
noise_pred = noise_pred_uncond + guidance_scale * noise_guidance
|
||||
noise_pred = noise_pred_uncond + guidance_scale * noise_guidance
|
||||
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
# compute the previous noisy sample x_t -> x_t-1
|
||||
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample
|
||||
|
||||
# call the callback, if provided
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
# call the callback, if provided
|
||||
if (i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0:
|
||||
progress_bar.update()
|
||||
if callback is not None and i % callback_steps == 0:
|
||||
callback(i, t, latents)
|
||||
|
||||
# 8. Post-processing
|
||||
image = self.decode_latents(latents)
|
||||
|
||||
@@ -28,7 +28,9 @@ def get_down_block(
|
||||
resnet_groups=None,
|
||||
cross_attention_dim=None,
|
||||
downsample_padding=None,
|
||||
dual_cross_attention=None,
|
||||
dual_cross_attention=False,
|
||||
use_linear_projection=False,
|
||||
only_cross_attention=False,
|
||||
):
|
||||
down_block_type = down_block_type[7:] if down_block_type.startswith("UNetRes") else down_block_type
|
||||
if down_block_type == "DownBlockFlat":
|
||||
@@ -58,6 +60,9 @@ def get_down_block(
|
||||
downsample_padding=downsample_padding,
|
||||
cross_attention_dim=cross_attention_dim,
|
||||
attn_num_head_channels=attn_num_head_channels,
|
||||
dual_cross_attention=dual_cross_attention,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention,
|
||||
)
|
||||
raise ValueError(f"{down_block_type} is not supported.")
|
||||
|
||||
@@ -75,7 +80,9 @@ def get_up_block(
|
||||
attn_num_head_channels,
|
||||
resnet_groups=None,
|
||||
cross_attention_dim=None,
|
||||
dual_cross_attention=None,
|
||||
dual_cross_attention=False,
|
||||
use_linear_projection=False,
|
||||
only_cross_attention=False,
|
||||
):
|
||||
up_block_type = up_block_type[7:] if up_block_type.startswith("UNetRes") else up_block_type
|
||||
if up_block_type == "UpBlockFlat":
|
||||
@@ -105,6 +112,9 @@ def get_up_block(
|
||||
resnet_groups=resnet_groups,
|
||||
cross_attention_dim=cross_attention_dim,
|
||||
attn_num_head_channels=attn_num_head_channels,
|
||||
dual_cross_attention=dual_cross_attention,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention,
|
||||
)
|
||||
raise ValueError(f"{up_block_type} is not supported.")
|
||||
|
||||
@@ -166,6 +176,7 @@ class UNetFlatConditionModel(ModelMixin, ConfigMixin):
|
||||
"CrossAttnUpBlockFlat",
|
||||
"CrossAttnUpBlockFlat",
|
||||
),
|
||||
only_cross_attention: Union[bool, Tuple[bool]] = False,
|
||||
block_out_channels: Tuple[int] = (320, 640, 1280, 1280),
|
||||
layers_per_block: int = 2,
|
||||
downsample_padding: int = 1,
|
||||
@@ -177,6 +188,7 @@ class UNetFlatConditionModel(ModelMixin, ConfigMixin):
|
||||
attention_head_dim: Union[int, Tuple[int]] = 8,
|
||||
dual_cross_attention: bool = False,
|
||||
use_linear_projection: bool = False,
|
||||
num_class_embeds: Optional[int] = None,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -192,10 +204,17 @@ class UNetFlatConditionModel(ModelMixin, ConfigMixin):
|
||||
|
||||
self.time_embedding = TimestepEmbedding(timestep_input_dim, time_embed_dim)
|
||||
|
||||
# class embedding
|
||||
if num_class_embeds is not None:
|
||||
self.class_embedding = nn.Embedding(num_class_embeds, time_embed_dim)
|
||||
|
||||
self.down_blocks = nn.ModuleList([])
|
||||
self.mid_block = None
|
||||
self.up_blocks = nn.ModuleList([])
|
||||
|
||||
if isinstance(only_cross_attention, bool):
|
||||
only_cross_attention = [only_cross_attention] * len(down_block_types)
|
||||
|
||||
if isinstance(attention_head_dim, int):
|
||||
attention_head_dim = (attention_head_dim,) * len(down_block_types)
|
||||
|
||||
@@ -221,6 +240,7 @@ class UNetFlatConditionModel(ModelMixin, ConfigMixin):
|
||||
downsample_padding=downsample_padding,
|
||||
dual_cross_attention=dual_cross_attention,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention[i],
|
||||
)
|
||||
self.down_blocks.append(down_block)
|
||||
|
||||
@@ -245,6 +265,7 @@ class UNetFlatConditionModel(ModelMixin, ConfigMixin):
|
||||
# up
|
||||
reversed_block_out_channels = list(reversed(block_out_channels))
|
||||
reversed_attention_head_dim = list(reversed(attention_head_dim))
|
||||
only_cross_attention = list(reversed(only_cross_attention))
|
||||
output_channel = reversed_block_out_channels[0]
|
||||
for i, up_block_type in enumerate(up_block_types):
|
||||
is_final_block = i == len(block_out_channels) - 1
|
||||
@@ -275,6 +296,7 @@ class UNetFlatConditionModel(ModelMixin, ConfigMixin):
|
||||
attn_num_head_channels=reversed_attention_head_dim[i],
|
||||
dual_cross_attention=dual_cross_attention,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention[i],
|
||||
)
|
||||
self.up_blocks.append(up_block)
|
||||
prev_output_channel = output_channel
|
||||
@@ -285,15 +307,17 @@ class UNetFlatConditionModel(ModelMixin, ConfigMixin):
|
||||
self.conv_out = LinearMultiDim(block_out_channels[0], out_channels, kernel_size=3, padding=1)
|
||||
|
||||
def set_attention_slice(self, slice_size):
|
||||
if slice_size is not None and self.config.attention_head_dim % slice_size != 0:
|
||||
head_dims = self.config.attention_head_dim
|
||||
head_dims = [head_dims] if isinstance(head_dims, int) else head_dims
|
||||
if slice_size is not None and any(dim % slice_size != 0 for dim in head_dims):
|
||||
raise ValueError(
|
||||
f"Make sure slice_size {slice_size} is a divisor of "
|
||||
f"the number of heads used in cross_attention {self.config.attention_head_dim}"
|
||||
f"Make sure slice_size {slice_size} is a common divisor of "
|
||||
f"the number of heads used in cross_attention: {head_dims}"
|
||||
)
|
||||
if slice_size is not None and slice_size > self.config.attention_head_dim:
|
||||
if slice_size is not None and slice_size > min(head_dims):
|
||||
raise ValueError(
|
||||
f"Chunk_size {slice_size} has to be smaller or equal to "
|
||||
f"the number of heads used in cross_attention {self.config.attention_head_dim}"
|
||||
f"slice_size {slice_size} has to be smaller or equal to "
|
||||
f"the lowest number of heads used in cross_attention: min({head_dims}) = {min(head_dims)}"
|
||||
)
|
||||
|
||||
for block in self.down_blocks:
|
||||
@@ -326,6 +350,7 @@ class UNetFlatConditionModel(ModelMixin, ConfigMixin):
|
||||
sample: torch.FloatTensor,
|
||||
timestep: Union[torch.Tensor, float, int],
|
||||
encoder_hidden_states: torch.Tensor,
|
||||
class_labels: Optional[torch.Tensor] = None,
|
||||
return_dict: bool = True,
|
||||
) -> Union[UNet2DConditionOutput, Tuple]:
|
||||
r"""
|
||||
@@ -378,6 +403,12 @@ class UNetFlatConditionModel(ModelMixin, ConfigMixin):
|
||||
t_emb = t_emb.to(dtype=self.dtype)
|
||||
emb = self.time_embedding(t_emb)
|
||||
|
||||
if self.config.num_class_embeds is not None:
|
||||
if class_labels is None:
|
||||
raise ValueError("class_labels should be provided when num_class_embeds > 0")
|
||||
class_emb = self.class_embedding(class_labels).to(dtype=self.dtype)
|
||||
emb = emb + class_emb
|
||||
|
||||
# 2. pre-process
|
||||
sample = self.conv_in(sample)
|
||||
|
||||
@@ -648,6 +679,7 @@ class CrossAttnDownBlockFlat(nn.Module):
|
||||
add_downsample=True,
|
||||
dual_cross_attention=False,
|
||||
use_linear_projection=False,
|
||||
only_cross_attention=False,
|
||||
):
|
||||
super().__init__()
|
||||
resnets = []
|
||||
@@ -682,6 +714,7 @@ class CrossAttnDownBlockFlat(nn.Module):
|
||||
cross_attention_dim=cross_attention_dim,
|
||||
norm_num_groups=resnet_groups,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention,
|
||||
)
|
||||
)
|
||||
else:
|
||||
@@ -712,15 +745,17 @@ class CrossAttnDownBlockFlat(nn.Module):
|
||||
self.gradient_checkpointing = False
|
||||
|
||||
def set_attention_slice(self, slice_size):
|
||||
if slice_size is not None and self.attn_num_head_channels % slice_size != 0:
|
||||
head_dims = self.attn_num_head_channels
|
||||
head_dims = [head_dims] if isinstance(head_dims, int) else head_dims
|
||||
if slice_size is not None and any(dim % slice_size != 0 for dim in head_dims):
|
||||
raise ValueError(
|
||||
f"Make sure slice_size {slice_size} is a divisor of "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"Make sure slice_size {slice_size} is a common divisor of "
|
||||
f"the number of heads used in cross_attention: {head_dims}"
|
||||
)
|
||||
if slice_size is not None and slice_size > self.attn_num_head_channels:
|
||||
if slice_size is not None and slice_size > min(head_dims):
|
||||
raise ValueError(
|
||||
f"Chunk_size {slice_size} has to be smaller or equal to "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"slice_size {slice_size} has to be smaller or equal to "
|
||||
f"the lowest number of heads used in cross_attention: min({head_dims}) = {min(head_dims)}"
|
||||
)
|
||||
|
||||
for attn in self.attentions:
|
||||
@@ -861,6 +896,7 @@ class CrossAttnUpBlockFlat(nn.Module):
|
||||
add_upsample=True,
|
||||
dual_cross_attention=False,
|
||||
use_linear_projection=False,
|
||||
only_cross_attention=False,
|
||||
):
|
||||
super().__init__()
|
||||
resnets = []
|
||||
@@ -897,6 +933,7 @@ class CrossAttnUpBlockFlat(nn.Module):
|
||||
cross_attention_dim=cross_attention_dim,
|
||||
norm_num_groups=resnet_groups,
|
||||
use_linear_projection=use_linear_projection,
|
||||
only_cross_attention=only_cross_attention,
|
||||
)
|
||||
)
|
||||
else:
|
||||
@@ -921,15 +958,17 @@ class CrossAttnUpBlockFlat(nn.Module):
|
||||
self.gradient_checkpointing = False
|
||||
|
||||
def set_attention_slice(self, slice_size):
|
||||
if slice_size is not None and self.attn_num_head_channels % slice_size != 0:
|
||||
head_dims = self.attn_num_head_channels
|
||||
head_dims = [head_dims] if isinstance(head_dims, int) else head_dims
|
||||
if slice_size is not None and any(dim % slice_size != 0 for dim in head_dims):
|
||||
raise ValueError(
|
||||
f"Make sure slice_size {slice_size} is a divisor of "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"Make sure slice_size {slice_size} is a common divisor of "
|
||||
f"the number of heads used in cross_attention: {head_dims}"
|
||||
)
|
||||
if slice_size is not None and slice_size > self.attn_num_head_channels:
|
||||
if slice_size is not None and slice_size > min(head_dims):
|
||||
raise ValueError(
|
||||
f"Chunk_size {slice_size} has to be smaller or equal to "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"slice_size {slice_size} has to be smaller or equal to "
|
||||
f"the lowest number of heads used in cross_attention: min({head_dims}) = {min(head_dims)}"
|
||||
)
|
||||
|
||||
for attn in self.attentions:
|
||||
@@ -1000,7 +1039,6 @@ class UNetMidBlockFlatCrossAttn(nn.Module):
|
||||
cross_attention_dim=1280,
|
||||
dual_cross_attention=False,
|
||||
use_linear_projection=False,
|
||||
**kwargs,
|
||||
):
|
||||
super().__init__()
|
||||
|
||||
@@ -1068,15 +1106,17 @@ class UNetMidBlockFlatCrossAttn(nn.Module):
|
||||
self.resnets = nn.ModuleList(resnets)
|
||||
|
||||
def set_attention_slice(self, slice_size):
|
||||
if slice_size is not None and self.attn_num_head_channels % slice_size != 0:
|
||||
head_dims = self.attn_num_head_channels
|
||||
head_dims = [head_dims] if isinstance(head_dims, int) else head_dims
|
||||
if slice_size is not None and any(dim % slice_size != 0 for dim in head_dims):
|
||||
raise ValueError(
|
||||
f"Make sure slice_size {slice_size} is a divisor of "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"Make sure slice_size {slice_size} is a common divisor of "
|
||||
f"the number of heads used in cross_attention: {head_dims}"
|
||||
)
|
||||
if slice_size is not None and slice_size > self.attn_num_head_channels:
|
||||
if slice_size is not None and slice_size > min(head_dims):
|
||||
raise ValueError(
|
||||
f"Chunk_size {slice_size} has to be smaller or equal to "
|
||||
f"the number of heads used in cross_attention {self.attn_num_head_channels}"
|
||||
f"slice_size {slice_size} has to be smaller or equal to "
|
||||
f"the lowest number of heads used in cross_attention: min({head_dims}) = {min(head_dims)}"
|
||||
)
|
||||
|
||||
for attn in self.attentions:
|
||||
|
||||
@@ -78,6 +78,7 @@ class VersatileDiffusionPipeline(DiffusionPipeline):
|
||||
vae=vae,
|
||||
scheduler=scheduler,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
|
||||
def enable_attention_slicing(self, slice_size: Optional[Union[str, int]] = "auto"):
|
||||
r"""
|
||||
@@ -111,8 +112,8 @@ class VersatileDiffusionPipeline(DiffusionPipeline):
|
||||
def image_variation(
|
||||
self,
|
||||
image: Union[torch.FloatTensor, PIL.Image.Image],
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -131,9 +132,9 @@ class VersatileDiffusionPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
image (`PIL.Image.Image`, `List[PIL.Image.Image]` or `torch.Tensor`):
|
||||
The image prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -193,7 +194,7 @@ class VersatileDiffusionPipeline(DiffusionPipeline):
|
||||
>>> pipe = pipe.to("cuda")
|
||||
|
||||
>>> generator = torch.Generator(device="cuda").manual_seed(0)
|
||||
>>> image = pipe(image, generator=generator).images[0]
|
||||
>>> image = pipe.image_variation(image, generator=generator).images[0]
|
||||
>>> image.save("./car_variation.png")
|
||||
```
|
||||
|
||||
@@ -227,8 +228,8 @@ class VersatileDiffusionPipeline(DiffusionPipeline):
|
||||
def text_to_image(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -247,9 +248,9 @@ class VersatileDiffusionPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -341,8 +342,8 @@ class VersatileDiffusionPipeline(DiffusionPipeline):
|
||||
prompt: Union[PIL.Image.Image, List[PIL.Image.Image]],
|
||||
image: Union[str, List[str]],
|
||||
text_to_image_strength: float = 0.5,
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
num_images_per_prompt: Optional[int] = 1,
|
||||
@@ -360,9 +361,9 @@ class VersatileDiffusionPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
|
||||
@@ -65,6 +65,8 @@ class VersatileDiffusionDualGuidedPipeline(DiffusionPipeline):
|
||||
vae: AutoencoderKL
|
||||
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler]
|
||||
|
||||
_optional_components = ["text_unet"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
tokenizer: CLIPTokenizer,
|
||||
@@ -87,6 +89,7 @@ class VersatileDiffusionDualGuidedPipeline(DiffusionPipeline):
|
||||
vae=vae,
|
||||
scheduler=scheduler,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
|
||||
if self.text_unet is not None and (
|
||||
"dual_cross_attention" not in self.image_unet.config or not self.image_unet.config.dual_cross_attention
|
||||
@@ -142,6 +145,8 @@ class VersatileDiffusionDualGuidedPipeline(DiffusionPipeline):
|
||||
index = int(index)
|
||||
self.image_unet.get_submodule(parent_name)[index] = module.transformers[0]
|
||||
|
||||
self.image_unet.register_to_config(dual_cross_attention=False)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_xformers_memory_efficient_attention with unet->image_unet
|
||||
def enable_xformers_memory_efficient_attention(self):
|
||||
r"""
|
||||
@@ -177,9 +182,14 @@ class VersatileDiffusionDualGuidedPipeline(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.image_unet.config.attention_head_dim // 2
|
||||
if isinstance(self.image_unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.image_unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.image_unet.config.attention_head_dim)
|
||||
|
||||
self.image_unet.set_attention_slice(slice_size)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_attention_slicing
|
||||
@@ -419,7 +429,7 @@ class VersatileDiffusionDualGuidedPipeline(DiffusionPipeline):
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_latents
|
||||
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
|
||||
shape = (batch_size, num_channels_latents, height // 8, width // 8)
|
||||
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor)
|
||||
if latents is None:
|
||||
if device.type == "mps":
|
||||
# randn does not work reproducibly on mps
|
||||
@@ -454,8 +464,8 @@ class VersatileDiffusionDualGuidedPipeline(DiffusionPipeline):
|
||||
prompt: Union[PIL.Image.Image, List[PIL.Image.Image]],
|
||||
image: Union[str, List[str]],
|
||||
text_to_image_strength: float = 0.5,
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
num_images_per_prompt: Optional[int] = 1,
|
||||
@@ -474,9 +484,9 @@ class VersatileDiffusionDualGuidedPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -551,6 +561,9 @@ class VersatileDiffusionDualGuidedPipeline(DiffusionPipeline):
|
||||
[`~pipelines.stable_diffusion.ImagePipelineOutput`] if `return_dict` is True, otherwise a `tuple. When
|
||||
returning a tuple, the first element is a list with the generated images.
|
||||
"""
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.image_unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.image_unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
# 1. Check inputs. Raise error if not correct
|
||||
self.check_inputs(prompt, image, height, width, callback_steps)
|
||||
|
||||
@@ -71,6 +71,7 @@ class VersatileDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
vae=vae,
|
||||
scheduler=scheduler,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_xformers_memory_efficient_attention with unet->image_unet
|
||||
def enable_xformers_memory_efficient_attention(self):
|
||||
@@ -107,9 +108,14 @@ class VersatileDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.image_unet.config.attention_head_dim // 2
|
||||
if isinstance(self.image_unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.image_unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.image_unet.config.attention_head_dim)
|
||||
|
||||
self.image_unet.set_attention_slice(slice_size)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_attention_slicing
|
||||
@@ -277,7 +283,7 @@ class VersatileDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_latents
|
||||
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
|
||||
shape = (batch_size, num_channels_latents, height // 8, width // 8)
|
||||
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor)
|
||||
if latents is None:
|
||||
if device.type == "mps":
|
||||
# randn does not work reproducibly on mps
|
||||
@@ -297,8 +303,8 @@ class VersatileDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
image: Union[PIL.Image.Image, List[PIL.Image.Image], torch.Tensor],
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -318,9 +324,9 @@ class VersatileDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
Args:
|
||||
image (`PIL.Image.Image`, `List[PIL.Image.Image]` or `torch.Tensor`):
|
||||
The image prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -391,6 +397,9 @@ class VersatileDiffusionImageVariationPipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.image_unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.image_unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
# 1. Check inputs. Raise error if not correct
|
||||
self.check_inputs(image, height, width, callback_steps)
|
||||
|
||||
@@ -57,6 +57,8 @@ class VersatileDiffusionTextToImagePipeline(DiffusionPipeline):
|
||||
vae: AutoencoderKL
|
||||
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler]
|
||||
|
||||
_optional_components = ["text_unet"]
|
||||
|
||||
def __init__(
|
||||
self,
|
||||
tokenizer: CLIPTokenizer,
|
||||
@@ -75,6 +77,7 @@ class VersatileDiffusionTextToImagePipeline(DiffusionPipeline):
|
||||
vae=vae,
|
||||
scheduler=scheduler,
|
||||
)
|
||||
self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1)
|
||||
|
||||
if self.text_unet is not None:
|
||||
self._swap_unet_attention_blocks()
|
||||
@@ -130,9 +133,14 @@ class VersatileDiffusionTextToImagePipeline(DiffusionPipeline):
|
||||
`attention_head_dim` must be a multiple of `slice_size`.
|
||||
"""
|
||||
if slice_size == "auto":
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.image_unet.config.attention_head_dim // 2
|
||||
if isinstance(self.image_unet.config.attention_head_dim, int):
|
||||
# half the attention head size is usually a good trade-off between
|
||||
# speed and memory
|
||||
slice_size = self.image_unet.config.attention_head_dim // 2
|
||||
else:
|
||||
# if `attention_head_dim` is a list, take the smallest head size
|
||||
slice_size = min(self.image_unet.config.attention_head_dim)
|
||||
|
||||
self.image_unet.set_attention_slice(slice_size)
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_attention_slicing
|
||||
@@ -337,7 +345,7 @@ class VersatileDiffusionTextToImagePipeline(DiffusionPipeline):
|
||||
|
||||
# Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_latents
|
||||
def prepare_latents(self, batch_size, num_channels_latents, height, width, dtype, device, generator, latents=None):
|
||||
shape = (batch_size, num_channels_latents, height // 8, width // 8)
|
||||
shape = (batch_size, num_channels_latents, height // self.vae_scale_factor, width // self.vae_scale_factor)
|
||||
if latents is None:
|
||||
if device.type == "mps":
|
||||
# randn does not work reproducibly on mps
|
||||
@@ -357,8 +365,8 @@ class VersatileDiffusionTextToImagePipeline(DiffusionPipeline):
|
||||
def __call__(
|
||||
self,
|
||||
prompt: Union[str, List[str]],
|
||||
height: int = 512,
|
||||
width: int = 512,
|
||||
height: Optional[int] = None,
|
||||
width: Optional[int] = None,
|
||||
num_inference_steps: int = 50,
|
||||
guidance_scale: float = 7.5,
|
||||
negative_prompt: Optional[Union[str, List[str]]] = None,
|
||||
@@ -378,9 +386,9 @@ class VersatileDiffusionTextToImagePipeline(DiffusionPipeline):
|
||||
Args:
|
||||
prompt (`str` or `List[str]`):
|
||||
The prompt or prompts to guide the image generation.
|
||||
height (`int`, *optional*, defaults to 512):
|
||||
height (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The height in pixels of the generated image.
|
||||
width (`int`, *optional*, defaults to 512):
|
||||
width (`int`, *optional*, defaults to self.image_unet.config.sample_size * self.vae_scale_factor):
|
||||
The width in pixels of the generated image.
|
||||
num_inference_steps (`int`, *optional*, defaults to 50):
|
||||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||||
@@ -443,6 +451,9 @@ class VersatileDiffusionTextToImagePipeline(DiffusionPipeline):
|
||||
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
|
||||
(nsfw) content, according to the `safety_checker`.
|
||||
"""
|
||||
# 0. Default height and width to unet
|
||||
height = height or self.image_unet.config.sample_size * self.vae_scale_factor
|
||||
width = width or self.image_unet.config.sample_size * self.vae_scale_factor
|
||||
|
||||
# 1. Check inputs. Raise error if not correct
|
||||
self.check_inputs(prompt, height, width, callback_steps)
|
||||
|
||||
@@ -22,6 +22,7 @@ if is_torch_available():
|
||||
from .scheduling_dpmsolver_multistep import DPMSolverMultistepScheduler
|
||||
from .scheduling_euler_ancestral_discrete import EulerAncestralDiscreteScheduler
|
||||
from .scheduling_euler_discrete import EulerDiscreteScheduler
|
||||
from .scheduling_heun import HeunDiscreteScheduler
|
||||
from .scheduling_ipndm import IPNDMScheduler
|
||||
from .scheduling_karras_ve import KarrasVeScheduler
|
||||
from .scheduling_pndm import PNDMScheduler
|
||||
|
||||
@@ -17,13 +17,13 @@
|
||||
|
||||
import math
|
||||
from dataclasses import dataclass
|
||||
from typing import Optional, Tuple, Union
|
||||
from typing import List, Optional, Tuple, Union
|
||||
|
||||
import numpy as np
|
||||
import torch
|
||||
|
||||
from ..configuration_utils import ConfigMixin, register_to_config
|
||||
from ..utils import _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS, BaseOutput
|
||||
from ..utils import _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS, BaseOutput, deprecate
|
||||
from .scheduling_utils import SchedulerMixin
|
||||
|
||||
|
||||
@@ -106,10 +106,15 @@ class DDIMScheduler(SchedulerMixin, ConfigMixin):
|
||||
an offset added to the inference steps. You can use a combination of `offset=1` and
|
||||
`set_alpha_to_one=False`, to make the last step use step 0 for the previous alpha product, as done in
|
||||
stable diffusion.
|
||||
prediction_type (`str`, default `epsilon`):
|
||||
indicates whether the model predicts the noise (epsilon), or the samples. One of `epsilon`, `sample`.
|
||||
`v-prediction` is not supported for this scheduler.
|
||||
|
||||
"""
|
||||
|
||||
_compatibles = _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
_deprecated_kwargs = ["predict_epsilon"]
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
@@ -118,13 +123,23 @@ class DDIMScheduler(SchedulerMixin, ConfigMixin):
|
||||
beta_start: float = 0.0001,
|
||||
beta_end: float = 0.02,
|
||||
beta_schedule: str = "linear",
|
||||
trained_betas: Optional[np.ndarray] = None,
|
||||
trained_betas: Optional[Union[np.ndarray, List[float]]] = None,
|
||||
clip_sample: bool = True,
|
||||
set_alpha_to_one: bool = True,
|
||||
steps_offset: int = 0,
|
||||
prediction_type: str = "epsilon",
|
||||
**kwargs,
|
||||
):
|
||||
message = (
|
||||
"Please make sure to instantiate your scheduler with `prediction_type` instead. E.g. `scheduler ="
|
||||
" DDIMScheduler.from_pretrained(<model_id>, prediction_type='epsilon')`."
|
||||
)
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.11.0", message, take_from=kwargs)
|
||||
if predict_epsilon is not None:
|
||||
self.register_to_config(prediction_type="epsilon" if predict_epsilon else "sample")
|
||||
|
||||
if trained_betas is not None:
|
||||
self.betas = torch.from_numpy(trained_betas)
|
||||
self.betas = torch.tensor(trained_betas, dtype=torch.float32)
|
||||
elif beta_schedule == "linear":
|
||||
self.betas = torch.linspace(beta_start, beta_end, num_train_timesteps, dtype=torch.float32)
|
||||
elif beta_schedule == "scaled_linear":
|
||||
@@ -258,7 +273,19 @@ class DDIMScheduler(SchedulerMixin, ConfigMixin):
|
||||
|
||||
# 3. compute predicted original sample from predicted noise also called
|
||||
# "predicted x_0" of formula (12) from https://arxiv.org/pdf/2010.02502.pdf
|
||||
pred_original_sample = (sample - beta_prod_t ** (0.5) * model_output) / alpha_prod_t ** (0.5)
|
||||
if self.config.prediction_type == "epsilon":
|
||||
pred_original_sample = (sample - beta_prod_t ** (0.5) * model_output) / alpha_prod_t ** (0.5)
|
||||
elif self.config.prediction_type == "sample":
|
||||
pred_original_sample = model_output
|
||||
elif self.config.prediction_type == "v_prediction":
|
||||
pred_original_sample = (alpha_prod_t**0.5) * sample - (beta_prod_t**0.5) * model_output
|
||||
# predict V
|
||||
model_output = (alpha_prod_t**0.5) * model_output + (beta_prod_t**0.5) * sample
|
||||
else:
|
||||
raise ValueError(
|
||||
f"prediction_type given as {self.config.prediction_type} must be one of `epsilon`, `sample`, or"
|
||||
" `v_prediction`"
|
||||
)
|
||||
|
||||
# 4. Clip "predicted x_0"
|
||||
if self.config.clip_sample:
|
||||
@@ -329,5 +356,25 @@ class DDIMScheduler(SchedulerMixin, ConfigMixin):
|
||||
noisy_samples = sqrt_alpha_prod * original_samples + sqrt_one_minus_alpha_prod * noise
|
||||
return noisy_samples
|
||||
|
||||
def get_velocity(
|
||||
self, sample: torch.FloatTensor, noise: torch.FloatTensor, timesteps: torch.IntTensor
|
||||
) -> torch.FloatTensor:
|
||||
# Make sure alphas_cumprod and timestep have same device and dtype as sample
|
||||
self.alphas_cumprod = self.alphas_cumprod.to(device=sample.device, dtype=sample.dtype)
|
||||
timesteps = timesteps.to(sample.device)
|
||||
|
||||
sqrt_alpha_prod = self.alphas_cumprod[timesteps] ** 0.5
|
||||
sqrt_alpha_prod = sqrt_alpha_prod.flatten()
|
||||
while len(sqrt_alpha_prod.shape) < len(sample.shape):
|
||||
sqrt_alpha_prod = sqrt_alpha_prod.unsqueeze(-1)
|
||||
|
||||
sqrt_one_minus_alpha_prod = (1 - self.alphas_cumprod[timesteps]) ** 0.5
|
||||
sqrt_one_minus_alpha_prod = sqrt_one_minus_alpha_prod.flatten()
|
||||
while len(sqrt_one_minus_alpha_prod.shape) < len(sample.shape):
|
||||
sqrt_one_minus_alpha_prod = sqrt_one_minus_alpha_prod.unsqueeze(-1)
|
||||
|
||||
velocity = sqrt_alpha_prod * noise - sqrt_one_minus_alpha_prod * sample
|
||||
return velocity
|
||||
|
||||
def __len__(self):
|
||||
return self.config.num_train_timesteps
|
||||
|
||||
@@ -23,6 +23,7 @@ import flax
|
||||
import jax.numpy as jnp
|
||||
|
||||
from ..configuration_utils import ConfigMixin, register_to_config
|
||||
from ..utils import deprecate
|
||||
from .scheduling_utils_flax import (
|
||||
_FLAX_COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS,
|
||||
FlaxSchedulerMixin,
|
||||
@@ -108,9 +109,14 @@ class FlaxDDIMScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
an offset added to the inference steps. You can use a combination of `offset=1` and
|
||||
`set_alpha_to_one=False`, to make the last step use step 0 for the previous alpha product, as done in
|
||||
stable diffusion.
|
||||
prediction_type (`str`, default `epsilon`):
|
||||
indicates whether the model predicts the noise (epsilon), or the samples. One of `epsilon`, `sample`.
|
||||
`v-prediction` is not supported for this scheduler.
|
||||
|
||||
"""
|
||||
|
||||
_compatibles = _FLAX_COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
_deprecated_kwargs = ["predict_epsilon"]
|
||||
|
||||
@property
|
||||
def has_state(self):
|
||||
@@ -125,7 +131,17 @@ class FlaxDDIMScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
beta_schedule: str = "linear",
|
||||
set_alpha_to_one: bool = True,
|
||||
steps_offset: int = 0,
|
||||
prediction_type: str = "epsilon",
|
||||
**kwargs,
|
||||
):
|
||||
message = (
|
||||
"Please make sure to instantiate your scheduler with `prediction_type` instead. E.g. `scheduler ="
|
||||
" FlaxDDIMScheduler.from_pretrained(<model_id>, prediction_type='epsilon')`."
|
||||
)
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.11.0", message, take_from=kwargs)
|
||||
if predict_epsilon is not None:
|
||||
self.register_to_config(prediction_type="epsilon" if predict_epsilon else "sample")
|
||||
|
||||
if beta_schedule == "linear":
|
||||
self.betas = jnp.linspace(beta_start, beta_end, num_train_timesteps, dtype=jnp.float32)
|
||||
elif beta_schedule == "scaled_linear":
|
||||
@@ -259,7 +275,19 @@ class FlaxDDIMScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
|
||||
# 3. compute predicted original sample from predicted noise also called
|
||||
# "predicted x_0" of formula (12) from https://arxiv.org/pdf/2010.02502.pdf
|
||||
pred_original_sample = (sample - beta_prod_t ** (0.5) * model_output) / alpha_prod_t ** (0.5)
|
||||
if self.config.prediction_type == "epsilon":
|
||||
pred_original_sample = (sample - beta_prod_t ** (0.5) * model_output) / alpha_prod_t ** (0.5)
|
||||
elif self.config.prediction_type == "sample":
|
||||
pred_original_sample = model_output
|
||||
elif self.config.prediction_type == "v_prediction":
|
||||
pred_original_sample = (alpha_prod_t**0.5) * sample - (beta_prod_t**0.5) * model_output
|
||||
# predict V
|
||||
model_output = (alpha_prod_t**0.5) * model_output + (beta_prod_t**0.5) * sample
|
||||
else:
|
||||
raise ValueError(
|
||||
f"prediction_type given as {self.config.prediction_type} must be one of `epsilon`, `sample`, or"
|
||||
" `v_prediction`"
|
||||
)
|
||||
|
||||
# 4. compute variance: "sigma_t(η)" -> see formula (16)
|
||||
# σ_t = sqrt((1 − α_t−1)/(1 − α_t)) * sqrt(1 − α_t/α_t−1)
|
||||
|
||||
@@ -16,7 +16,7 @@
|
||||
|
||||
import math
|
||||
from dataclasses import dataclass
|
||||
from typing import Optional, Tuple, Union
|
||||
from typing import List, Optional, Tuple, Union
|
||||
|
||||
import numpy as np
|
||||
import torch
|
||||
@@ -99,12 +99,14 @@ class DDPMScheduler(SchedulerMixin, ConfigMixin):
|
||||
`fixed_small_log`, `fixed_large`, `fixed_large_log`, `learned` or `learned_range`.
|
||||
clip_sample (`bool`, default `True`):
|
||||
option to clip predicted sample between -1 and 1 for numerical stability.
|
||||
predict_epsilon (`bool`):
|
||||
optional flag to use when the model predicts the noise (epsilon), or the samples instead of the noise.
|
||||
|
||||
prediction_type (`str`, default `epsilon`):
|
||||
indicates whether the model predicts the noise (epsilon), or the samples. One of `epsilon`, `sample`.
|
||||
`v-prediction` is not supported for this scheduler.
|
||||
"""
|
||||
|
||||
_compatibles = _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
_deprecated_kwargs = ["predict_epsilon"]
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
@@ -113,13 +115,22 @@ class DDPMScheduler(SchedulerMixin, ConfigMixin):
|
||||
beta_start: float = 0.0001,
|
||||
beta_end: float = 0.02,
|
||||
beta_schedule: str = "linear",
|
||||
trained_betas: Optional[np.ndarray] = None,
|
||||
trained_betas: Optional[Union[np.ndarray, List[float]]] = None,
|
||||
variance_type: str = "fixed_small",
|
||||
clip_sample: bool = True,
|
||||
predict_epsilon: bool = True,
|
||||
prediction_type: str = "epsilon",
|
||||
**kwargs,
|
||||
):
|
||||
message = (
|
||||
"Please make sure to instantiate your scheduler with `prediction_type` instead. E.g. `scheduler ="
|
||||
" DDPMScheduler.from_pretrained(<model_id>, prediction_type='epsilon')`."
|
||||
)
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.11.0", message, take_from=kwargs)
|
||||
if predict_epsilon is not None:
|
||||
self.register_to_config(prediction_type="epsilon" if predict_epsilon else "sample")
|
||||
|
||||
if trained_betas is not None:
|
||||
self.betas = torch.from_numpy(trained_betas)
|
||||
self.betas = torch.tensor(trained_betas, dtype=torch.float32)
|
||||
elif beta_schedule == "linear":
|
||||
self.betas = torch.linspace(beta_start, beta_end, num_train_timesteps, dtype=torch.float32)
|
||||
elif beta_schedule == "scaled_linear":
|
||||
@@ -241,13 +252,13 @@ class DDPMScheduler(SchedulerMixin, ConfigMixin):
|
||||
|
||||
"""
|
||||
message = (
|
||||
"Please make sure to instantiate your scheduler with `predict_epsilon` instead. E.g. `scheduler ="
|
||||
" DDPMScheduler.from_pretrained(<model_id>, predict_epsilon=True)`."
|
||||
"Please make sure to instantiate your scheduler with `prediction_type` instead. E.g. `scheduler ="
|
||||
" DDPMScheduler.from_pretrained(<model_id>, prediction_type='epsilon')`."
|
||||
)
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.10.0", message, take_from=kwargs)
|
||||
if predict_epsilon is not None and predict_epsilon != self.config.predict_epsilon:
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.11.0", message, take_from=kwargs)
|
||||
if predict_epsilon is not None:
|
||||
new_config = dict(self.config)
|
||||
new_config["predict_epsilon"] = predict_epsilon
|
||||
new_config["prediction_type"] = "epsilon" if predict_epsilon else "sample"
|
||||
self._internal_dict = FrozenDict(new_config)
|
||||
|
||||
t = timestep
|
||||
@@ -265,10 +276,15 @@ class DDPMScheduler(SchedulerMixin, ConfigMixin):
|
||||
|
||||
# 2. compute predicted original sample from predicted noise also called
|
||||
# "predicted x_0" of formula (15) from https://arxiv.org/pdf/2006.11239.pdf
|
||||
if self.config.predict_epsilon:
|
||||
if self.config.prediction_type == "epsilon":
|
||||
pred_original_sample = (sample - beta_prod_t ** (0.5) * model_output) / alpha_prod_t ** (0.5)
|
||||
else:
|
||||
elif self.config.prediction_type == "sample":
|
||||
pred_original_sample = model_output
|
||||
else:
|
||||
raise ValueError(
|
||||
f"prediction_type given as {self.config.prediction_type} must be one of `epsilon`, `sample` "
|
||||
" for the DDPMScheduler."
|
||||
)
|
||||
|
||||
# 3. Clip "predicted x_0"
|
||||
if self.config.clip_sample:
|
||||
@@ -330,5 +346,25 @@ class DDPMScheduler(SchedulerMixin, ConfigMixin):
|
||||
noisy_samples = sqrt_alpha_prod * original_samples + sqrt_one_minus_alpha_prod * noise
|
||||
return noisy_samples
|
||||
|
||||
def get_velocity(
|
||||
self, sample: torch.FloatTensor, noise: torch.FloatTensor, timesteps: torch.IntTensor
|
||||
) -> torch.FloatTensor:
|
||||
# Make sure alphas_cumprod and timestep have same device and dtype as sample
|
||||
self.alphas_cumprod = self.alphas_cumprod.to(device=sample.device, dtype=sample.dtype)
|
||||
timesteps = timesteps.to(sample.device)
|
||||
|
||||
sqrt_alpha_prod = self.alphas_cumprod[timesteps] ** 0.5
|
||||
sqrt_alpha_prod = sqrt_alpha_prod.flatten()
|
||||
while len(sqrt_alpha_prod.shape) < len(sample.shape):
|
||||
sqrt_alpha_prod = sqrt_alpha_prod.unsqueeze(-1)
|
||||
|
||||
sqrt_one_minus_alpha_prod = (1 - self.alphas_cumprod[timesteps]) ** 0.5
|
||||
sqrt_one_minus_alpha_prod = sqrt_one_minus_alpha_prod.flatten()
|
||||
while len(sqrt_one_minus_alpha_prod.shape) < len(sample.shape):
|
||||
sqrt_one_minus_alpha_prod = sqrt_one_minus_alpha_prod.unsqueeze(-1)
|
||||
|
||||
velocity = sqrt_alpha_prod * noise - sqrt_one_minus_alpha_prod * sample
|
||||
return velocity
|
||||
|
||||
def __len__(self):
|
||||
return self.config.num_train_timesteps
|
||||
|
||||
@@ -103,12 +103,13 @@ class FlaxDDPMScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
`fixed_small_log`, `fixed_large`, `fixed_large_log`, `learned` or `learned_range`.
|
||||
clip_sample (`bool`, default `True`):
|
||||
option to clip predicted sample between -1 and 1 for numerical stability.
|
||||
predict_epsilon (`bool`):
|
||||
optional flag to use when the model predicts the noise (epsilon), or the samples instead of the noise.
|
||||
|
||||
prediction_type (`str`, default `epsilon`):
|
||||
indicates whether the model predicts the noise (epsilon), or the samples. One of `epsilon`, `sample`.
|
||||
`v-prediction` is not supported for this scheduler.
|
||||
"""
|
||||
|
||||
_compatibles = _FLAX_COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
_deprecated_kwargs = ["predict_epsilon"]
|
||||
|
||||
@property
|
||||
def has_state(self):
|
||||
@@ -124,8 +125,17 @@ class FlaxDDPMScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
trained_betas: Optional[jnp.ndarray] = None,
|
||||
variance_type: str = "fixed_small",
|
||||
clip_sample: bool = True,
|
||||
predict_epsilon: bool = True,
|
||||
prediction_type: str = "epsilon",
|
||||
**kwargs,
|
||||
):
|
||||
message = (
|
||||
"Please make sure to instantiate your scheduler with `prediction_type` instead. E.g. `scheduler ="
|
||||
" FlaxDDPMScheduler.from_pretrained(<model_id>, prediction_type='epsilon')`."
|
||||
)
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.11.0", message, take_from=kwargs)
|
||||
if predict_epsilon is not None:
|
||||
self.register_to_config(prediction_type="epsilon" if predict_epsilon else "sample")
|
||||
|
||||
if trained_betas is not None:
|
||||
self.betas = jnp.asarray(trained_betas)
|
||||
elif beta_schedule == "linear":
|
||||
@@ -204,7 +214,6 @@ class FlaxDDPMScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
timestep: int,
|
||||
sample: jnp.ndarray,
|
||||
key: random.KeyArray,
|
||||
predict_epsilon: bool = True,
|
||||
return_dict: bool = True,
|
||||
**kwargs,
|
||||
) -> Union[FlaxDDPMSchedulerOutput, Tuple]:
|
||||
@@ -227,13 +236,13 @@ class FlaxDDPMScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
|
||||
"""
|
||||
message = (
|
||||
"Please make sure to instantiate your scheduler with `predict_epsilon` instead. E.g. `scheduler ="
|
||||
" DDPMScheduler.from_pretrained(<model_id>, predict_epsilon=True)`."
|
||||
"Please make sure to instantiate your scheduler with `prediction_type` instead. E.g. `scheduler ="
|
||||
" FlaxDDPMScheduler.from_pretrained(<model_id>, prediction_type='epsilon')`."
|
||||
)
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.10.0", message, take_from=kwargs)
|
||||
if predict_epsilon is not None and predict_epsilon != self.config.predict_epsilon:
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.11.0", message, take_from=kwargs)
|
||||
if predict_epsilon is not None:
|
||||
new_config = dict(self.config)
|
||||
new_config["predict_epsilon"] = predict_epsilon
|
||||
new_config["prediction_type"] = "epsilon" if predict_epsilon else "sample"
|
||||
self._internal_dict = FrozenDict(new_config)
|
||||
|
||||
t = timestep
|
||||
@@ -251,10 +260,15 @@ class FlaxDDPMScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
|
||||
# 2. compute predicted original sample from predicted noise also called
|
||||
# "predicted x_0" of formula (15) from https://arxiv.org/pdf/2006.11239.pdf
|
||||
if self.config.predict_epsilon:
|
||||
if self.config.prediction_type == "epsilon":
|
||||
pred_original_sample = (sample - beta_prod_t ** (0.5) * model_output) / alpha_prod_t ** (0.5)
|
||||
else:
|
||||
elif self.config.prediction_type == "sample":
|
||||
pred_original_sample = model_output
|
||||
else:
|
||||
raise ValueError(
|
||||
f"prediction_type given as {self.config.prediction_type} must be one of `epsilon`, `sample` "
|
||||
" for the FlaxDDPMScheduler."
|
||||
)
|
||||
|
||||
# 3. Clip "predicted x_0"
|
||||
if self.config.clip_sample:
|
||||
|
||||
@@ -21,7 +21,7 @@ import numpy as np
|
||||
import torch
|
||||
|
||||
from ..configuration_utils import ConfigMixin, register_to_config
|
||||
from ..utils import _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS
|
||||
from ..utils import _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS, deprecate
|
||||
from .scheduling_utils import SchedulerMixin, SchedulerOutput
|
||||
|
||||
|
||||
@@ -87,10 +87,9 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
|
||||
solver_order (`int`, default `2`):
|
||||
the order of DPM-Solver; can be `1` or `2` or `3`. We recommend to use `solver_order=2` for guided
|
||||
sampling, and `solver_order=3` for unconditional sampling.
|
||||
predict_epsilon (`bool`, default `True`):
|
||||
we currently support both the noise prediction model and the data prediction model. If the model predicts
|
||||
the noise / epsilon, set `predict_epsilon` to `True`. If the model predicts the data / x0 directly, set
|
||||
`predict_epsilon` to `False`.
|
||||
prediction_type (`str`, default `epsilon`):
|
||||
indicates whether the model predicts the noise (epsilon), or the data / `x0`. One of `epsilon`, `sample`,
|
||||
or `v-prediction`.
|
||||
thresholding (`bool`, default `False`):
|
||||
whether to use the "dynamic thresholding" method (introduced by Imagen, https://arxiv.org/abs/2205.11487).
|
||||
For pixel-space diffusion models, you can set both `algorithm_type=dpmsolver++` and `thresholding=True` to
|
||||
@@ -118,6 +117,8 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
|
||||
_compatibles = _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
_deprecated_kwargs = ["predict_epsilon"]
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
@@ -126,18 +127,27 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
|
||||
beta_start: float = 0.0001,
|
||||
beta_end: float = 0.02,
|
||||
beta_schedule: str = "linear",
|
||||
trained_betas: Optional[np.ndarray] = None,
|
||||
trained_betas: Optional[Union[np.ndarray, List[float]]] = None,
|
||||
solver_order: int = 2,
|
||||
predict_epsilon: bool = True,
|
||||
prediction_type: str = "epsilon",
|
||||
thresholding: bool = False,
|
||||
dynamic_thresholding_ratio: float = 0.995,
|
||||
sample_max_value: float = 1.0,
|
||||
algorithm_type: str = "dpmsolver++",
|
||||
solver_type: str = "midpoint",
|
||||
lower_order_final: bool = True,
|
||||
**kwargs,
|
||||
):
|
||||
message = (
|
||||
"Please make sure to instantiate your scheduler with `prediction_type` instead. E.g. `scheduler ="
|
||||
" DPMSolverMultistepScheduler.from_pretrained(<model_id>, prediction_type='epsilon')`."
|
||||
)
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.11.0", message, take_from=kwargs)
|
||||
if predict_epsilon is not None:
|
||||
self.register_to_config(prediction_type="epsilon" if predict_epsilon else "sample")
|
||||
|
||||
if trained_betas is not None:
|
||||
self.betas = torch.from_numpy(trained_betas)
|
||||
self.betas = torch.tensor(trained_betas, dtype=torch.float32)
|
||||
elif beta_schedule == "linear":
|
||||
self.betas = torch.linspace(beta_start, beta_end, num_train_timesteps, dtype=torch.float32)
|
||||
elif beta_schedule == "scaled_linear":
|
||||
@@ -203,7 +213,7 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
Convert the model output to the corresponding type that the algorithm (DPM-Solver / DPM-Solver++) needs.
|
||||
|
||||
DPM-Solver is designed to discretize an integral of the noise prediciton model, and DPM-Solver++ is designed to
|
||||
DPM-Solver is designed to discretize an integral of the noise prediction model, and DPM-Solver++ is designed to
|
||||
discretize an integral of the data prediction model. So we need to first convert the model output to the
|
||||
corresponding type to match the algorithm.
|
||||
|
||||
@@ -221,13 +231,25 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
# DPM-Solver++ needs to solve an integral of the data prediction model.
|
||||
if self.config.algorithm_type == "dpmsolver++":
|
||||
if self.config.predict_epsilon:
|
||||
if self.config.prediction_type == "epsilon":
|
||||
alpha_t, sigma_t = self.alpha_t[timestep], self.sigma_t[timestep]
|
||||
x0_pred = (sample - sigma_t * model_output) / alpha_t
|
||||
else:
|
||||
elif self.config.prediction_type == "sample":
|
||||
x0_pred = model_output
|
||||
elif self.config.prediction_type == "v_prediction":
|
||||
alpha_t, sigma_t = self.alpha_t[timestep], self.sigma_t[timestep]
|
||||
x0_pred = alpha_t * sample - sigma_t * model_output
|
||||
else:
|
||||
raise ValueError(
|
||||
f"prediction_type given as {self.config.prediction_type} must be one of `epsilon`, `sample`, or"
|
||||
" `v_prediction` for the DPMSolverMultistepScheduler."
|
||||
)
|
||||
|
||||
if self.config.thresholding:
|
||||
# Dynamic thresholding in https://arxiv.org/abs/2205.11487
|
||||
orig_dtype = x0_pred.dtype
|
||||
if orig_dtype not in [torch.float, torch.double]:
|
||||
x0_pred = x0_pred.float()
|
||||
dynamic_max_val = torch.quantile(
|
||||
torch.abs(x0_pred).reshape((x0_pred.shape[0], -1)), self.config.dynamic_thresholding_ratio, dim=1
|
||||
)
|
||||
@@ -236,15 +258,25 @@ class DPMSolverMultistepScheduler(SchedulerMixin, ConfigMixin):
|
||||
self.config.sample_max_value * torch.ones_like(dynamic_max_val).to(dynamic_max_val.device),
|
||||
)[(...,) + (None,) * (x0_pred.ndim - 1)]
|
||||
x0_pred = torch.clamp(x0_pred, -dynamic_max_val, dynamic_max_val) / dynamic_max_val
|
||||
x0_pred = x0_pred.type(orig_dtype)
|
||||
return x0_pred
|
||||
# DPM-Solver needs to solve an integral of the noise prediction model.
|
||||
elif self.config.algorithm_type == "dpmsolver":
|
||||
if self.config.predict_epsilon:
|
||||
if self.config.prediction_type == "epsilon":
|
||||
return model_output
|
||||
else:
|
||||
elif self.config.prediction_type == "sample":
|
||||
alpha_t, sigma_t = self.alpha_t[timestep], self.sigma_t[timestep]
|
||||
epsilon = (sample - alpha_t * model_output) / sigma_t
|
||||
return epsilon
|
||||
elif self.config.prediction_type == "v_prediction":
|
||||
alpha_t, sigma_t = self.alpha_t[timestep], self.sigma_t[timestep]
|
||||
epsilon = alpha_t * model_output + sigma_t * sample
|
||||
return epsilon
|
||||
else:
|
||||
raise ValueError(
|
||||
f"prediction_type given as {self.config.prediction_type} must be one of `epsilon`, `sample`, or"
|
||||
" `v_prediction` for the DPMSolverMultistepScheduler."
|
||||
)
|
||||
|
||||
def dpm_solver_first_order_update(
|
||||
self,
|
||||
|
||||
@@ -23,6 +23,7 @@ import jax
|
||||
import jax.numpy as jnp
|
||||
|
||||
from ..configuration_utils import ConfigMixin, register_to_config
|
||||
from ..utils import deprecate
|
||||
from .scheduling_utils_flax import (
|
||||
_FLAX_COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS,
|
||||
FlaxSchedulerMixin,
|
||||
@@ -118,10 +119,9 @@ class FlaxDPMSolverMultistepScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
solver_order (`int`, default `2`):
|
||||
the order of DPM-Solver; can be `1` or `2` or `3`. We recommend to use `solver_order=2` for guided
|
||||
sampling, and `solver_order=3` for unconditional sampling.
|
||||
predict_epsilon (`bool`, default `True`):
|
||||
we currently support both the noise prediction model and the data prediction model. If the model predicts
|
||||
the noise / epsilon, set `predict_epsilon` to `True`. If the model predicts the data / x0 directly, set
|
||||
`predict_epsilon` to `False`.
|
||||
prediction_type (`str`, default `epsilon`):
|
||||
indicates whether the model predicts the noise (epsilon), or the data / `x0`. One of `epsilon`, `sample`,
|
||||
or `v-prediction`.
|
||||
thresholding (`bool`, default `False`):
|
||||
whether to use the "dynamic thresholding" method (introduced by Imagen, https://arxiv.org/abs/2205.11487).
|
||||
For pixel-space diffusion models, you can set both `algorithm_type=dpmsolver++` and `thresholding=True` to
|
||||
@@ -149,6 +149,7 @@ class FlaxDPMSolverMultistepScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
|
||||
_compatibles = _FLAX_COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
_deprecated_kwargs = ["predict_epsilon"]
|
||||
|
||||
@property
|
||||
def has_state(self):
|
||||
@@ -163,14 +164,23 @@ class FlaxDPMSolverMultistepScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
beta_schedule: str = "linear",
|
||||
trained_betas: Optional[jnp.ndarray] = None,
|
||||
solver_order: int = 2,
|
||||
predict_epsilon: bool = True,
|
||||
prediction_type: str = "epsilon",
|
||||
thresholding: bool = False,
|
||||
dynamic_thresholding_ratio: float = 0.995,
|
||||
sample_max_value: float = 1.0,
|
||||
algorithm_type: str = "dpmsolver++",
|
||||
solver_type: str = "midpoint",
|
||||
lower_order_final: bool = True,
|
||||
**kwargs,
|
||||
):
|
||||
message = (
|
||||
"Please make sure to instantiate your scheduler with `prediction_type` instead. E.g. `scheduler ="
|
||||
" FlaxDPMSolverMultistepScheduler.from_pretrained(<model_id>, prediction_type='epsilon')`."
|
||||
)
|
||||
predict_epsilon = deprecate("predict_epsilon", "0.11.0", message, take_from=kwargs)
|
||||
if predict_epsilon is not None:
|
||||
self.register_to_config(prediction_type="epsilon" if predict_epsilon else "sample")
|
||||
|
||||
if trained_betas is not None:
|
||||
self.betas = jnp.asarray(trained_betas)
|
||||
elif beta_schedule == "linear":
|
||||
@@ -242,7 +252,7 @@ class FlaxDPMSolverMultistepScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
Convert the model output to the corresponding type that the algorithm (DPM-Solver / DPM-Solver++) needs.
|
||||
|
||||
DPM-Solver is designed to discretize an integral of the noise prediciton model, and DPM-Solver++ is designed to
|
||||
DPM-Solver is designed to discretize an integral of the noise prediction model, and DPM-Solver++ is designed to
|
||||
discretize an integral of the data prediction model. So we need to first convert the model output to the
|
||||
corresponding type to match the algorithm.
|
||||
|
||||
@@ -260,11 +270,20 @@ class FlaxDPMSolverMultistepScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
# DPM-Solver++ needs to solve an integral of the data prediction model.
|
||||
if self.config.algorithm_type == "dpmsolver++":
|
||||
if self.config.predict_epsilon:
|
||||
if self.config.prediction_type == "epsilon":
|
||||
alpha_t, sigma_t = self.alpha_t[timestep], self.sigma_t[timestep]
|
||||
x0_pred = (sample - sigma_t * model_output) / alpha_t
|
||||
else:
|
||||
elif self.config.prediction_type == "sample":
|
||||
x0_pred = model_output
|
||||
elif self.config.prediction_type == "v_prediction":
|
||||
alpha_t, sigma_t = self.alpha_t[timestep], self.sigma_t[timestep]
|
||||
x0_pred = alpha_t * sample - sigma_t * model_output
|
||||
else:
|
||||
raise ValueError(
|
||||
f"prediction_type given as {self.config.prediction_type} must be one of `epsilon`, `sample`, "
|
||||
" or `v_prediction` for the FlaxDPMSolverMultistepScheduler."
|
||||
)
|
||||
|
||||
if self.config.thresholding:
|
||||
# Dynamic thresholding in https://arxiv.org/abs/2205.11487
|
||||
dynamic_max_val = jnp.percentile(
|
||||
@@ -277,12 +296,21 @@ class FlaxDPMSolverMultistepScheduler(FlaxSchedulerMixin, ConfigMixin):
|
||||
return x0_pred
|
||||
# DPM-Solver needs to solve an integral of the noise prediction model.
|
||||
elif self.config.algorithm_type == "dpmsolver":
|
||||
if self.config.predict_epsilon:
|
||||
if self.config.prediction_type == "epsilon":
|
||||
return model_output
|
||||
else:
|
||||
elif self.config.prediction_type == "sample":
|
||||
alpha_t, sigma_t = self.alpha_t[timestep], self.sigma_t[timestep]
|
||||
epsilon = (sample - alpha_t * model_output) / sigma_t
|
||||
return epsilon
|
||||
elif self.config.prediction_type == "v_prediction":
|
||||
alpha_t, sigma_t = self.alpha_t[timestep], self.sigma_t[timestep]
|
||||
epsilon = alpha_t * model_output + sigma_t * sample
|
||||
return epsilon
|
||||
else:
|
||||
raise ValueError(
|
||||
f"prediction_type given as {self.config.prediction_type} must be one of `epsilon`, `sample`, "
|
||||
" or `v_prediction` for the FlaxDPMSolverMultistepScheduler."
|
||||
)
|
||||
|
||||
def dpm_solver_first_order_update(
|
||||
self, model_output: jnp.ndarray, timestep: int, prev_timestep: int, sample: jnp.ndarray
|
||||
|
||||
@@ -13,7 +13,7 @@
|
||||
# limitations under the License.
|
||||
|
||||
from dataclasses import dataclass
|
||||
from typing import Optional, Tuple, Union
|
||||
from typing import List, Optional, Tuple, Union
|
||||
|
||||
import numpy as np
|
||||
import torch
|
||||
@@ -68,6 +68,7 @@ class EulerAncestralDiscreteScheduler(SchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
|
||||
_compatibles = _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
@@ -76,10 +77,10 @@ class EulerAncestralDiscreteScheduler(SchedulerMixin, ConfigMixin):
|
||||
beta_start: float = 0.0001,
|
||||
beta_end: float = 0.02,
|
||||
beta_schedule: str = "linear",
|
||||
trained_betas: Optional[np.ndarray] = None,
|
||||
trained_betas: Optional[Union[np.ndarray, List[float]]] = None,
|
||||
):
|
||||
if trained_betas is not None:
|
||||
self.betas = torch.from_numpy(trained_betas)
|
||||
self.betas = torch.tensor(trained_betas, dtype=torch.float32)
|
||||
elif beta_schedule == "linear":
|
||||
self.betas = torch.linspace(beta_start, beta_end, num_train_timesteps, dtype=torch.float32)
|
||||
elif beta_schedule == "scaled_linear":
|
||||
|
||||
@@ -13,7 +13,7 @@
|
||||
# limitations under the License.
|
||||
|
||||
from dataclasses import dataclass
|
||||
from typing import Optional, Tuple, Union
|
||||
from typing import List, Optional, Tuple, Union
|
||||
|
||||
import numpy as np
|
||||
import torch
|
||||
@@ -69,6 +69,7 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
|
||||
_compatibles = _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
@@ -77,10 +78,11 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
|
||||
beta_start: float = 0.0001,
|
||||
beta_end: float = 0.02,
|
||||
beta_schedule: str = "linear",
|
||||
trained_betas: Optional[np.ndarray] = None,
|
||||
trained_betas: Optional[Union[np.ndarray, List[float]]] = None,
|
||||
prediction_type: str = "epsilon",
|
||||
):
|
||||
if trained_betas is not None:
|
||||
self.betas = torch.from_numpy(trained_betas)
|
||||
self.betas = torch.tensor(trained_betas, dtype=torch.float32)
|
||||
elif beta_schedule == "linear":
|
||||
self.betas = torch.linspace(beta_start, beta_end, num_train_timesteps, dtype=torch.float32)
|
||||
elif beta_schedule == "scaled_linear":
|
||||
@@ -229,7 +231,15 @@ class EulerDiscreteScheduler(SchedulerMixin, ConfigMixin):
|
||||
sample = sample + eps * (sigma_hat**2 - sigma**2) ** 0.5
|
||||
|
||||
# 1. compute predicted original sample (x_0) from sigma-scaled predicted noise
|
||||
pred_original_sample = sample - sigma_hat * model_output
|
||||
if self.config.prediction_type == "epsilon":
|
||||
pred_original_sample = sample - sigma_hat * model_output
|
||||
elif self.config.prediction_type == "v_prediction":
|
||||
# * c_out + input * c_skip
|
||||
pred_original_sample = model_output * (-sigma / (sigma**2 + 1) ** 0.5) + (sample / (sigma**2 + 1))
|
||||
else:
|
||||
raise ValueError(
|
||||
f"prediction_type given as {self.config.prediction_type} must be one of `epsilon`, or `v_prediction`"
|
||||
)
|
||||
|
||||
# 2. Convert to an ODE derivative
|
||||
derivative = (sample - pred_original_sample) / sigma_hat
|
||||
|
||||
247
src/diffusers/schedulers/scheduling_heun.py
Normal file
247
src/diffusers/schedulers/scheduling_heun.py
Normal file
@@ -0,0 +1,247 @@
|
||||
# Copyright 2022 Katherine Crowson, The HuggingFace Team and hlky. All rights reserved.
|
||||
#
|
||||
# Licensed under the Apache License, Version 2.0 (the "License");
|
||||
# you may not use this file except in compliance with the License.
|
||||
# You may obtain a copy of the License at
|
||||
#
|
||||
# http://www.apache.org/licenses/LICENSE-2.0
|
||||
#
|
||||
# Unless required by applicable law or agreed to in writing, software
|
||||
# distributed under the License is distributed on an "AS IS" BASIS,
|
||||
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
|
||||
# See the License for the specific language governing permissions and
|
||||
# limitations under the License.
|
||||
|
||||
from typing import List, Optional, Tuple, Union
|
||||
|
||||
import numpy as np
|
||||
import torch
|
||||
|
||||
from ..configuration_utils import ConfigMixin, register_to_config
|
||||
from ..utils import _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS
|
||||
from .scheduling_utils import SchedulerMixin, SchedulerOutput
|
||||
|
||||
|
||||
class HeunDiscreteScheduler(SchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
Args:
|
||||
Implements Algorithm 2 (Heun steps) from Karras et al. (2022). for discrete beta schedules. Based on the original
|
||||
k-diffusion implementation by Katherine Crowson:
|
||||
https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L90
|
||||
[`~ConfigMixin`] takes care of storing all config attributes that are passed in the scheduler's `__init__`
|
||||
function, such as `num_train_timesteps`. They can be accessed via `scheduler.config.num_train_timesteps`.
|
||||
[`~ConfigMixin`] also provides general loading and saving functionality via the [`~ConfigMixin.save_config`] and
|
||||
[`~ConfigMixin.from_config`] functions.
|
||||
num_train_timesteps (`int`): number of diffusion steps used to train the model. beta_start (`float`): the
|
||||
starting `beta` value of inference. beta_end (`float`): the final `beta` value. beta_schedule (`str`):
|
||||
the beta schedule, a mapping from a beta range to a sequence of betas for stepping the model. Choose from
|
||||
`linear` or `scaled_linear`.
|
||||
trained_betas (`np.ndarray`, optional):
|
||||
option to pass an array of betas directly to the constructor to bypass `beta_start`, `beta_end` etc.
|
||||
options to clip the variance used when adding noise to the denoised sample. Choose from `fixed_small`,
|
||||
`fixed_small_log`, `fixed_large`, `fixed_large_log`, `learned` or `learned_range`.
|
||||
tensor_format (`str`): whether the scheduler expects pytorch or numpy arrays.
|
||||
"""
|
||||
|
||||
_compatibles = _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
order = 2
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
self,
|
||||
num_train_timesteps: int = 1000,
|
||||
beta_start: float = 0.00085, # sensible defaults
|
||||
beta_end: float = 0.012,
|
||||
beta_schedule: str = "linear",
|
||||
trained_betas: Optional[Union[np.ndarray, List[float]]] = None,
|
||||
):
|
||||
if trained_betas is not None:
|
||||
self.betas = torch.tensor(trained_betas, dtype=torch.float32)
|
||||
elif beta_schedule == "linear":
|
||||
self.betas = torch.linspace(beta_start, beta_end, num_train_timesteps, dtype=torch.float32)
|
||||
elif beta_schedule == "scaled_linear":
|
||||
# this schedule is very specific to the latent diffusion model.
|
||||
self.betas = (
|
||||
torch.linspace(beta_start**0.5, beta_end**0.5, num_train_timesteps, dtype=torch.float32) ** 2
|
||||
)
|
||||
else:
|
||||
raise NotImplementedError(f"{beta_schedule} does is not implemented for {self.__class__}")
|
||||
|
||||
self.alphas = 1.0 - self.betas
|
||||
self.alphas_cumprod = torch.cumprod(self.alphas, dim=0)
|
||||
|
||||
# set all values
|
||||
self.set_timesteps(num_train_timesteps, None, num_train_timesteps)
|
||||
|
||||
def index_for_timestep(self, timestep):
|
||||
indices = (self.timesteps == timestep).nonzero()
|
||||
if self.state_in_first_order:
|
||||
pos = 0 if indices.shape[0] < 2 else 1
|
||||
else:
|
||||
pos = 0
|
||||
return indices[pos].item()
|
||||
|
||||
def scale_model_input(
|
||||
self,
|
||||
sample: torch.FloatTensor,
|
||||
timestep: Union[float, torch.FloatTensor],
|
||||
) -> torch.FloatTensor:
|
||||
"""
|
||||
Args:
|
||||
Ensures interchangeability with schedulers that need to scale the denoising model input depending on the
|
||||
current timestep.
|
||||
sample (`torch.FloatTensor`): input sample timestep (`int`, optional): current timestep
|
||||
Returns:
|
||||
`torch.FloatTensor`: scaled input sample
|
||||
"""
|
||||
step_index = self.index_for_timestep(timestep)
|
||||
|
||||
sigma = self.sigmas[step_index]
|
||||
sample = sample / ((sigma**2 + 1) ** 0.5)
|
||||
return sample
|
||||
|
||||
def set_timesteps(
|
||||
self,
|
||||
num_inference_steps: int,
|
||||
device: Union[str, torch.device] = None,
|
||||
num_train_timesteps: Optional[int] = None,
|
||||
):
|
||||
"""
|
||||
Sets the timesteps used for the diffusion chain. Supporting function to be run before inference.
|
||||
|
||||
Args:
|
||||
num_inference_steps (`int`):
|
||||
the number of diffusion steps used when generating samples with a pre-trained model.
|
||||
device (`str` or `torch.device`, optional):
|
||||
the device to which the timesteps should be moved to. If `None`, the timesteps are not moved.
|
||||
"""
|
||||
self.num_inference_steps = num_inference_steps
|
||||
|
||||
num_train_timesteps = num_train_timesteps or self.config.num_train_timesteps
|
||||
|
||||
timesteps = np.linspace(0, num_train_timesteps - 1, num_inference_steps, dtype=float)[::-1].copy()
|
||||
|
||||
sigmas = np.array(((1 - self.alphas_cumprod) / self.alphas_cumprod) ** 0.5)
|
||||
sigmas = np.interp(timesteps, np.arange(0, len(sigmas)), sigmas)
|
||||
sigmas = np.concatenate([sigmas, [0.0]]).astype(np.float32)
|
||||
sigmas = torch.from_numpy(sigmas).to(device=device)
|
||||
self.sigmas = torch.cat([sigmas[:1], sigmas[1:-1].repeat_interleave(2), sigmas[-1:]])
|
||||
|
||||
# standard deviation of the initial noise distribution
|
||||
self.init_noise_sigma = self.sigmas.max()
|
||||
|
||||
timesteps = torch.from_numpy(timesteps)
|
||||
timesteps = torch.cat([timesteps[:1], timesteps[1:].repeat_interleave(2), timesteps[-1:]])
|
||||
|
||||
if str(device).startswith("mps"):
|
||||
# mps does not support float64
|
||||
self.timesteps = timesteps.to(device, dtype=torch.float32)
|
||||
else:
|
||||
self.timesteps = timesteps.to(device=device)
|
||||
|
||||
# empty dt and derivative
|
||||
self.prev_derivative = None
|
||||
self.dt = None
|
||||
|
||||
@property
|
||||
def state_in_first_order(self):
|
||||
return self.dt is None
|
||||
|
||||
def step(
|
||||
self,
|
||||
model_output: Union[torch.FloatTensor, np.ndarray],
|
||||
timestep: Union[float, torch.FloatTensor],
|
||||
sample: Union[torch.FloatTensor, np.ndarray],
|
||||
return_dict: bool = True,
|
||||
) -> Union[SchedulerOutput, Tuple]:
|
||||
"""
|
||||
Args:
|
||||
Predict the sample at the previous timestep by reversing the SDE. Core function to propagate the diffusion
|
||||
process from the learned model outputs (most often the predicted noise).
|
||||
model_output (`torch.FloatTensor` or `np.ndarray`): direct output from learned diffusion model. timestep
|
||||
(`int`): current discrete timestep in the diffusion chain. sample (`torch.FloatTensor` or `np.ndarray`):
|
||||
current instance of sample being created by diffusion process.
|
||||
return_dict (`bool`): option for returning tuple rather than SchedulerOutput class
|
||||
Returns:
|
||||
[`~schedulers.scheduling_utils.SchedulerOutput`] or `tuple`:
|
||||
[`~schedulers.scheduling_utils.SchedulerOutput`] if `return_dict` is True, otherwise a `tuple`. When
|
||||
returning a tuple, the first element is the sample tensor.
|
||||
"""
|
||||
step_index = self.index_for_timestep(timestep)
|
||||
|
||||
if self.state_in_first_order:
|
||||
sigma = self.sigmas[step_index]
|
||||
sigma_next = self.sigmas[step_index + 1]
|
||||
else:
|
||||
# 2nd order / Heun's method
|
||||
sigma = self.sigmas[step_index - 1]
|
||||
sigma_next = self.sigmas[step_index]
|
||||
|
||||
# currently only gamma=0 is supported. This usually works best anyways.
|
||||
# We can support gamma in the future but then need to scale the timestep before
|
||||
# passing it to the model which requires a change in API
|
||||
gamma = 0
|
||||
sigma_hat = sigma * (gamma + 1) # Note: sigma_hat == sigma for now
|
||||
|
||||
# 1. compute predicted original sample (x_0) from sigma-scaled predicted noise
|
||||
pred_original_sample = sample - sigma_hat * model_output
|
||||
|
||||
if self.state_in_first_order:
|
||||
# 2. Convert to an ODE derivative
|
||||
derivative = (sample - pred_original_sample) / sigma_hat
|
||||
# 3. 1st order derivative
|
||||
dt = sigma_next - sigma_hat
|
||||
|
||||
# store for 2nd order step
|
||||
self.prev_derivative = derivative
|
||||
self.dt = dt
|
||||
self.sample = sample
|
||||
else:
|
||||
# 2. 2nd order / Heun's method
|
||||
derivative = (sample - pred_original_sample) / sigma_hat
|
||||
derivative = (self.prev_derivative + derivative) / 2
|
||||
|
||||
# 3. Retrieve 1st order derivative
|
||||
dt = self.dt
|
||||
sample = self.sample
|
||||
|
||||
# free dt and derivative
|
||||
# Note, this puts the scheduler in "first order mode"
|
||||
self.prev_derivative = None
|
||||
self.dt = None
|
||||
self.sample = None
|
||||
|
||||
prev_sample = sample + derivative * dt
|
||||
|
||||
if not return_dict:
|
||||
return (prev_sample,)
|
||||
|
||||
return SchedulerOutput(prev_sample=prev_sample)
|
||||
|
||||
def add_noise(
|
||||
self,
|
||||
original_samples: torch.FloatTensor,
|
||||
noise: torch.FloatTensor,
|
||||
timesteps: torch.FloatTensor,
|
||||
) -> torch.FloatTensor:
|
||||
# Make sure sigmas and timesteps have the same device and dtype as original_samples
|
||||
self.sigmas = self.sigmas.to(device=original_samples.device, dtype=original_samples.dtype)
|
||||
if original_samples.device.type == "mps" and torch.is_floating_point(timesteps):
|
||||
# mps does not support float64
|
||||
self.timesteps = self.timesteps.to(original_samples.device, dtype=torch.float32)
|
||||
timesteps = timesteps.to(original_samples.device, dtype=torch.float32)
|
||||
else:
|
||||
self.timesteps = self.timesteps.to(original_samples.device)
|
||||
timesteps = timesteps.to(original_samples.device)
|
||||
|
||||
step_indices = [self.index_for_timestep(t) for t in timesteps]
|
||||
|
||||
sigma = self.sigmas[step_indices].flatten()
|
||||
while len(sigma.shape) < len(original_samples.shape):
|
||||
sigma = sigma.unsqueeze(-1)
|
||||
|
||||
noisy_samples = original_samples + noise * sigma
|
||||
return noisy_samples
|
||||
|
||||
def __len__(self):
|
||||
return self.config.num_train_timesteps
|
||||
@@ -13,8 +13,9 @@
|
||||
# limitations under the License.
|
||||
|
||||
import math
|
||||
from typing import Tuple, Union
|
||||
from typing import List, Optional, Tuple, Union
|
||||
|
||||
import numpy as np
|
||||
import torch
|
||||
|
||||
from ..configuration_utils import ConfigMixin, register_to_config
|
||||
@@ -37,8 +38,12 @@ class IPNDMScheduler(SchedulerMixin, ConfigMixin):
|
||||
num_train_timesteps (`int`): number of diffusion steps used to train the model.
|
||||
"""
|
||||
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(self, num_train_timesteps: int = 1000):
|
||||
def __init__(
|
||||
self, num_train_timesteps: int = 1000, trained_betas: Optional[Union[np.ndarray, List[float]]] = None
|
||||
):
|
||||
# set `betas`, `alphas`, `timesteps`
|
||||
self.set_timesteps(num_train_timesteps)
|
||||
|
||||
@@ -65,7 +70,11 @@ class IPNDMScheduler(SchedulerMixin, ConfigMixin):
|
||||
steps = torch.linspace(1, 0, num_inference_steps + 1)[:-1]
|
||||
steps = torch.cat([steps, torch.tensor([0.0])])
|
||||
|
||||
self.betas = torch.sin(steps * math.pi / 2) ** 2
|
||||
if self.config.trained_betas is not None:
|
||||
self.betas = torch.tensor(self.config.trained_betas, dtype=torch.float32)
|
||||
else:
|
||||
self.betas = torch.sin(steps * math.pi / 2) ** 2
|
||||
|
||||
self.alphas = (1.0 - self.betas**2) ** 0.5
|
||||
|
||||
timesteps = (torch.atan2(self.betas, self.alphas) / math.pi * 2)[:-1]
|
||||
|
||||
@@ -77,6 +77,8 @@ class KarrasVeScheduler(SchedulerMixin, ConfigMixin):
|
||||
|
||||
"""
|
||||
|
||||
order = 2
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
self,
|
||||
|
||||
@@ -13,7 +13,7 @@
|
||||
# limitations under the License.
|
||||
import warnings
|
||||
from dataclasses import dataclass
|
||||
from typing import Optional, Tuple, Union
|
||||
from typing import List, Optional, Tuple, Union
|
||||
|
||||
import numpy as np
|
||||
import torch
|
||||
@@ -68,6 +68,7 @@ class LMSDiscreteScheduler(SchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
|
||||
_compatibles = _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
@@ -76,10 +77,10 @@ class LMSDiscreteScheduler(SchedulerMixin, ConfigMixin):
|
||||
beta_start: float = 0.0001,
|
||||
beta_end: float = 0.02,
|
||||
beta_schedule: str = "linear",
|
||||
trained_betas: Optional[np.ndarray] = None,
|
||||
trained_betas: Optional[Union[np.ndarray, List[float]]] = None,
|
||||
):
|
||||
if trained_betas is not None:
|
||||
self.betas = torch.from_numpy(trained_betas)
|
||||
self.betas = torch.tensor(trained_betas, dtype=torch.float32)
|
||||
elif beta_schedule == "linear":
|
||||
self.betas = torch.linspace(beta_start, beta_end, num_train_timesteps, dtype=torch.float32)
|
||||
elif beta_schedule == "scaled_linear":
|
||||
|
||||
@@ -15,7 +15,7 @@
|
||||
# DISCLAIMER: This file is strongly influenced by https://github.com/ermongroup/ddim
|
||||
|
||||
import math
|
||||
from typing import Optional, Tuple, Union
|
||||
from typing import List, Optional, Tuple, Union
|
||||
|
||||
import numpy as np
|
||||
import torch
|
||||
@@ -90,6 +90,7 @@ class PNDMScheduler(SchedulerMixin, ConfigMixin):
|
||||
"""
|
||||
|
||||
_compatibles = _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS.copy()
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
@@ -98,13 +99,13 @@ class PNDMScheduler(SchedulerMixin, ConfigMixin):
|
||||
beta_start: float = 0.0001,
|
||||
beta_end: float = 0.02,
|
||||
beta_schedule: str = "linear",
|
||||
trained_betas: Optional[np.ndarray] = None,
|
||||
trained_betas: Optional[Union[np.ndarray, List[float]]] = None,
|
||||
skip_prk_steps: bool = False,
|
||||
set_alpha_to_one: bool = False,
|
||||
steps_offset: int = 0,
|
||||
):
|
||||
if trained_betas is not None:
|
||||
self.betas = torch.from_numpy(trained_betas)
|
||||
self.betas = torch.tensor(trained_betas, dtype=torch.float32)
|
||||
elif beta_schedule == "linear":
|
||||
self.betas = torch.linspace(beta_start, beta_end, num_train_timesteps, dtype=torch.float32)
|
||||
elif beta_schedule == "scaled_linear":
|
||||
|
||||
@@ -102,6 +102,8 @@ class RePaintScheduler(SchedulerMixin, ConfigMixin):
|
||||
|
||||
"""
|
||||
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
self,
|
||||
|
||||
@@ -66,6 +66,8 @@ class ScoreSdeVeScheduler(SchedulerMixin, ConfigMixin):
|
||||
correct_steps (`int`): number of correction steps performed on a produced sample.
|
||||
"""
|
||||
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
self,
|
||||
|
||||
@@ -38,6 +38,8 @@ class ScoreSdeVpScheduler(SchedulerMixin, ConfigMixin):
|
||||
|
||||
"""
|
||||
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(self, num_train_timesteps=2000, beta_min=0.1, beta_max=20, sampling_eps=1e-3):
|
||||
self.sigmas = None
|
||||
|
||||
@@ -138,6 +138,8 @@ class VQDiffusionScheduler(SchedulerMixin, ConfigMixin):
|
||||
The ending cumulative gamma value.
|
||||
"""
|
||||
|
||||
order = 1
|
||||
|
||||
@register_to_config
|
||||
def __init__(
|
||||
self,
|
||||
|
||||
@@ -28,6 +28,7 @@ from .import_utils import (
|
||||
is_inflect_available,
|
||||
is_modelcards_available,
|
||||
is_onnx_available,
|
||||
is_safetensors_available,
|
||||
is_scipy_available,
|
||||
is_tf_available,
|
||||
is_torch_available,
|
||||
@@ -69,6 +70,7 @@ CONFIG_NAME = "config.json"
|
||||
WEIGHTS_NAME = "diffusion_pytorch_model.bin"
|
||||
FLAX_WEIGHTS_NAME = "diffusion_flax_model.msgpack"
|
||||
ONNX_WEIGHTS_NAME = "model.onnx"
|
||||
SAFETENSORS_WEIGHTS_NAME = "diffusion_pytorch_model.safetensors"
|
||||
ONNX_EXTERNAL_WEIGHTS_NAME = "weights.pb"
|
||||
HUGGINGFACE_CO_RESOLVE_ENDPOINT = "https://huggingface.co"
|
||||
DIFFUSERS_CACHE = default_cache_path
|
||||
@@ -81,6 +83,7 @@ _COMPATIBLE_STABLE_DIFFUSION_SCHEDULERS = [
|
||||
"PNDMScheduler",
|
||||
"LMSDiscreteScheduler",
|
||||
"EulerDiscreteScheduler",
|
||||
"HeunDiscreteScheduler",
|
||||
"EulerAncestralDiscreteScheduler",
|
||||
"DPMSolverMultistepScheduler",
|
||||
]
|
||||
|
||||
@@ -32,7 +32,7 @@ def deprecate(*args, take_from: Optional[Union[Dict, Any]] = None, standard_warn
|
||||
|
||||
if warning is not None:
|
||||
warning = warning + " " if standard_warn else ""
|
||||
warnings.warn(warning + message, DeprecationWarning)
|
||||
warnings.warn(warning + message, FutureWarning, stacklevel=2)
|
||||
|
||||
if isinstance(deprecated_kwargs, dict) and len(deprecated_kwargs) > 0:
|
||||
call_frame = inspect.getouterframes(inspect.currentframe())[1]
|
||||
|
||||
@@ -362,6 +362,21 @@ class EulerDiscreteScheduler(metaclass=DummyObject):
|
||||
requires_backends(cls, ["torch"])
|
||||
|
||||
|
||||
class HeunDiscreteScheduler(metaclass=DummyObject):
|
||||
_backends = ["torch"]
|
||||
|
||||
def __init__(self, *args, **kwargs):
|
||||
requires_backends(self, ["torch"])
|
||||
|
||||
@classmethod
|
||||
def from_config(cls, *args, **kwargs):
|
||||
requires_backends(cls, ["torch"])
|
||||
|
||||
@classmethod
|
||||
def from_pretrained(cls, *args, **kwargs):
|
||||
requires_backends(cls, ["torch"])
|
||||
|
||||
|
||||
class IPNDMScheduler(metaclass=DummyObject):
|
||||
_backends = ["torch"]
|
||||
|
||||
|
||||
@@ -154,6 +154,21 @@ class StableDiffusionPipelineSafe(metaclass=DummyObject):
|
||||
requires_backends(cls, ["torch", "transformers"])
|
||||
|
||||
|
||||
class StableDiffusionUpscalePipeline(metaclass=DummyObject):
|
||||
_backends = ["torch", "transformers"]
|
||||
|
||||
def __init__(self, *args, **kwargs):
|
||||
requires_backends(self, ["torch", "transformers"])
|
||||
|
||||
@classmethod
|
||||
def from_config(cls, *args, **kwargs):
|
||||
requires_backends(cls, ["torch", "transformers"])
|
||||
|
||||
@classmethod
|
||||
def from_pretrained(cls, *args, **kwargs):
|
||||
requires_backends(cls, ["torch", "transformers"])
|
||||
|
||||
|
||||
class VersatileDiffusionDualGuidedPipeline(metaclass=DummyObject):
|
||||
_backends = ["torch", "transformers"]
|
||||
|
||||
|
||||
@@ -42,6 +42,7 @@ ENV_VARS_TRUE_AND_AUTO_VALUES = ENV_VARS_TRUE_VALUES.union({"AUTO"})
|
||||
USE_TF = os.environ.get("USE_TF", "AUTO").upper()
|
||||
USE_TORCH = os.environ.get("USE_TORCH", "AUTO").upper()
|
||||
USE_JAX = os.environ.get("USE_FLAX", "AUTO").upper()
|
||||
USE_SAFETENSORS = os.environ.get("USE_SAFETENSORS", "AUTO").upper()
|
||||
|
||||
STR_OPERATION_TO_FUNC = {">": op.gt, ">=": op.ge, "==": op.eq, "!=": op.ne, "<=": op.le, "<": op.lt}
|
||||
|
||||
@@ -55,7 +56,7 @@ if USE_TORCH in ENV_VARS_TRUE_AND_AUTO_VALUES and USE_TF not in ENV_VARS_TRUE_VA
|
||||
except importlib_metadata.PackageNotFoundError:
|
||||
_torch_available = False
|
||||
else:
|
||||
logger.info("Disabling PyTorch because USE_TF is set")
|
||||
logger.info("Disabling PyTorch because USE_TORCH is set")
|
||||
_torch_available = False
|
||||
|
||||
|
||||
@@ -109,6 +110,17 @@ if USE_JAX in ENV_VARS_TRUE_AND_AUTO_VALUES:
|
||||
else:
|
||||
_flax_available = False
|
||||
|
||||
if USE_SAFETENSORS in ENV_VARS_TRUE_AND_AUTO_VALUES:
|
||||
_safetensors_available = importlib.util.find_spec("safetensors") is not None
|
||||
if _safetensors_available:
|
||||
try:
|
||||
_safetensors_version = importlib_metadata.version("safetensors")
|
||||
logger.info(f"Safetensors version {_safetensors_version} available.")
|
||||
except importlib_metadata.PackageNotFoundError:
|
||||
_safetensors_available = False
|
||||
else:
|
||||
logger.info("Disabling Safetensors because USE_TF is set")
|
||||
_safetensors_available = False
|
||||
|
||||
_transformers_available = importlib.util.find_spec("transformers") is not None
|
||||
try:
|
||||
@@ -145,7 +157,13 @@ except importlib_metadata.PackageNotFoundError:
|
||||
_onnxruntime_version = "N/A"
|
||||
_onnx_available = importlib.util.find_spec("onnxruntime") is not None
|
||||
if _onnx_available:
|
||||
candidates = ("onnxruntime", "onnxruntime-gpu", "onnxruntime-directml", "onnxruntime-openvino")
|
||||
candidates = (
|
||||
"onnxruntime",
|
||||
"onnxruntime-gpu",
|
||||
"onnxruntime-directml",
|
||||
"onnxruntime-openvino",
|
||||
"ort_nightly_directml",
|
||||
)
|
||||
_onnxruntime_version = None
|
||||
# For the metadata, we have to look for both onnxruntime and onnxruntime-gpu
|
||||
for pkg in candidates:
|
||||
@@ -190,6 +208,10 @@ def is_torch_available():
|
||||
return _torch_available
|
||||
|
||||
|
||||
def is_safetensors_available():
|
||||
return _safetensors_available
|
||||
|
||||
|
||||
def is_tf_available():
|
||||
return _tf_available
|
||||
|
||||
|
||||
@@ -63,8 +63,8 @@ class UNet1DModelTests(ModelTesterMixin, unittest.TestCase):
|
||||
super().test_outputs_equivalence()
|
||||
|
||||
@unittest.skipIf(torch_device == "mps", "mish op not supported in MPS")
|
||||
def test_from_pretrained_save_pretrained(self):
|
||||
super().test_from_pretrained_save_pretrained()
|
||||
def test_from_save_pretrained(self):
|
||||
super().test_from_save_pretrained()
|
||||
|
||||
@unittest.skipIf(torch_device == "mps", "mish op not supported in MPS")
|
||||
def test_model_from_pretrained(self):
|
||||
@@ -183,8 +183,8 @@ class UNetRLModelTests(ModelTesterMixin, unittest.TestCase):
|
||||
super().test_outputs_equivalence()
|
||||
|
||||
@unittest.skipIf(torch_device == "mps", "mish op not supported in MPS")
|
||||
def test_from_pretrained_save_pretrained(self):
|
||||
super().test_from_pretrained_save_pretrained()
|
||||
def test_from_save_pretrained(self):
|
||||
super().test_from_save_pretrained()
|
||||
|
||||
@unittest.skipIf(torch_device == "mps", "mish op not supported in MPS")
|
||||
def test_model_from_pretrained(self):
|
||||
|
||||
@@ -639,3 +639,29 @@ class UNet2DConditionModelIntegrationTests(unittest.TestCase):
|
||||
expected_output_slice = torch.tensor(expected_slice)
|
||||
|
||||
assert torch_all_close(output_slice, expected_output_slice, atol=5e-3)
|
||||
|
||||
@parameterized.expand(
|
||||
[
|
||||
# fmt: off
|
||||
[83, 4, [0.1514, 0.0807, 0.1624, 0.1016, -0.1896, 0.0263, 0.0677, 0.2310]],
|
||||
[17, 0.55, [0.1164, -0.0216, 0.0170, 0.1589, -0.3120, 0.1005, -0.0581, -0.1458]],
|
||||
[8, 0.89, [-0.1758, -0.0169, 0.1004, -0.1411, 0.1312, 0.1103, -0.1996, 0.2139]],
|
||||
[3, 1000, [0.1214, 0.0352, -0.0731, -0.1562, -0.0994, -0.0906, -0.2340, -0.0539]],
|
||||
# fmt: on
|
||||
]
|
||||
)
|
||||
@require_torch_gpu
|
||||
def test_stabilityai_sd_v2_fp16(self, seed, timestep, expected_slice):
|
||||
model = self.get_unet_model(model_id="stabilityai/stable-diffusion-2", fp16=True)
|
||||
latents = self.get_latents(seed, shape=(4, 4, 96, 96), fp16=True)
|
||||
encoder_hidden_states = self.get_encoder_hidden_states(seed, shape=(4, 77, 1024), fp16=True)
|
||||
|
||||
with torch.no_grad():
|
||||
sample = model(latents, timestep=timestep, encoder_hidden_states=encoder_hidden_states).sample
|
||||
|
||||
assert sample.shape == latents.shape
|
||||
|
||||
output_slice = sample[-1, -2:, -2:, :2].flatten().float().cpu()
|
||||
expected_output_slice = torch.tensor(expected_slice)
|
||||
|
||||
assert torch_all_close(output_slice, expected_output_slice, atol=5e-3)
|
||||
|
||||
Some files were not shown because too many files have changed in this diff Show More
Reference in New Issue
Block a user