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2 Commits

Author SHA1 Message Date
Patrick von Platen
fb5ddb7fdc [Bad dependencies] Fix imports (#1382)
* fix imports

* better error

* up

* finish
2022-11-24 00:26:35 +01:00
anton-
2a49278a51 Patch release: v0.8.1 2022-11-24 00:12:23 +01:00
423 changed files with 11435 additions and 51801 deletions

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@@ -5,20 +5,7 @@ body:
- type: markdown
attributes:
value: |
Thanks a lot for taking the time to file this issue 🤗.
Issues do not only help to improve the library, but also publicly document common problems, questions, workflows for the whole community!
Thus, issues are of the same importance as pull requests when contributing to this library ❤️.
In order to make your issue as **useful for the community as possible**, let's try to stick to some simple guidelines:
- 1. Please try to be as precise and concise as possible.
*Give your issue a fitting title. Assume that someone which very limited knowledge of diffusers can understand your issue. Add links to the source code, documentation other issues, pull requests etc...*
- 2. If your issue is about something not working, **always** provide a reproducible code snippet. The reader should be able to reproduce your issue by **only copy-pasting your code snippet into a Python shell**.
*The community cannot solve your issue if it cannot reproduce it. If your bug is related to training, add your training script and make everything needed to train public. Otherwise, just add a simple Python code snippet.*
- 3. Add the **minimum amount of code / context that is needed to understand, reproduce your issue**.
*Make the life of maintainers easy. `diffusers` is getting many issues every day. Make sure your issue is about one bug and one bug only. Make sure you add only the context, code needed to understand your issues - nothing more. Generally, every issue is a way of documenting this library, try to make it a good documentation entry.*
- type: markdown
attributes:
value: |
For more in-detail information on how to write good issues you can have a look [here](https://huggingface.co/course/chapter8/5?fw=pt)
Thanks for taking the time to fill out this bug report!
- type: textarea
id: bug-description
attributes:
@@ -33,8 +20,6 @@ body:
label: Reproduction
description: Please provide a minimal reproducible code which we can copy/paste and reproduce the issue.
placeholder: Reproduction
validations:
required: true
- type: textarea
id: logs
attributes:

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@@ -13,6 +13,5 @@ jobs:
with:
commit_sha: ${{ github.sha }}
package: diffusers
languages: en ko
secrets:
token: ${{ secrets.HUGGINGFACE_PUSH }}

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@@ -14,4 +14,3 @@ jobs:
commit_sha: ${{ github.event.pull_request.head.sha }}
pr_number: ${{ github.event.number }}
package: diffusers
languages: en ko

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@@ -1,162 +0,0 @@
name: Nightly tests on main
on:
schedule:
- cron: "0 0 * * *" # every day at midnight
env:
DIFFUSERS_IS_CI: yes
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
RUN_NIGHTLY: yes
jobs:
run_nightly_tests:
strategy:
fail-fast: false
matrix:
config:
- name: Nightly PyTorch CUDA tests on Ubuntu
framework: pytorch
runner: docker-gpu
image: diffusers/diffusers-pytorch-cuda
report: torch_cuda
- name: Nightly Flax TPU tests on Ubuntu
framework: flax
runner: docker-tpu
image: diffusers/diffusers-flax-tpu
report: flax_tpu
- name: Nightly ONNXRuntime CUDA tests on Ubuntu
framework: onnxruntime
runner: docker-gpu
image: diffusers/diffusers-onnxruntime-cuda
report: onnx_cuda
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ ${{ matrix.config.runner == 'docker-tpu' && '--privileged' || '--gpus 0'}}
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
if: ${{ matrix.config.runner == 'docker-gpu' }}
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m pip install -e .[quality,test]
python -m pip install -U git+https://github.com/huggingface/transformers
python -m pip install git+https://github.com/huggingface/accelerate
- name: Environment
run: |
python utils/print_env.py
- name: Run nightly PyTorch CUDA tests
if: ${{ matrix.config.framework == 'pytorch' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Run nightly Flax TPU tests
if: ${{ matrix.config.framework == 'flax' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Run nightly ONNXRuntime CUDA tests
if: ${{ matrix.config.framework == 'onnxruntime' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: ${{ matrix.config.report }}_test_reports
path: reports
run_nightly_tests_apple_m1:
name: Nightly PyTorch MPS tests on MacOS
runs-on: [ self-hosted, apple-m1 ]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Clean checkout
shell: arch -arch arm64 bash {0}
run: |
git clean -fxd
- name: Setup miniconda
uses: ./.github/actions/setup-miniconda
with:
python-version: 3.9
- name: Install dependencies
shell: arch -arch arm64 bash {0}
run: |
${CONDA_RUN} python -m pip install --upgrade pip
${CONDA_RUN} python -m pip install -e .[quality,test]
${CONDA_RUN} python -m pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
${CONDA_RUN} python -m pip install git+https://github.com/huggingface/accelerate
- name: Environment
shell: arch -arch arm64 bash {0}
run: |
${CONDA_RUN} python utils/print_env.py
- name: Run nightly PyTorch tests on M1 (MPS)
shell: arch -arch arm64 bash {0}
env:
HF_HOME: /System/Volumes/Data/mnt/cache
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps tests/
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_torch_mps_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_mps_test_reports
path: reports

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@@ -1,4 +1,4 @@
name: Fast tests for PRs
name: Run fast tests
on:
pull_request:
@@ -14,6 +14,7 @@ env:
OMP_NUM_THREADS: 4
MKL_NUM_THREADS: 4
PYTEST_TIMEOUT: 60
MPS_TORCH_VERSION: 1.13.0
jobs:
run_fast_tests:
@@ -57,10 +58,9 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev -y
python -m pip install -e .[quality,test]
python -m pip install -U git+https://github.com/huggingface/transformers
python -m pip install git+https://github.com/huggingface/accelerate
python -m pip install -U git+https://github.com/huggingface/transformers
- name: Environment
run: |
@@ -126,7 +126,7 @@ jobs:
run: |
${CONDA_RUN} python -m pip install --upgrade pip
${CONDA_RUN} python -m pip install -e .[quality,test]
${CONDA_RUN} python -m pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
${CONDA_RUN} python -m pip install --pre torch==${MPS_TORCH_VERSION} --extra-index-url https://download.pytorch.org/whl/test/cpu
${CONDA_RUN} python -m pip install git+https://github.com/huggingface/accelerate
${CONDA_RUN} python -m pip install -U git+https://github.com/huggingface/transformers
@@ -137,9 +137,6 @@ jobs:
- name: Run fast PyTorch tests on M1 (MPS)
shell: arch -arch arm64 bash {0}
env:
HF_HOME: /System/Volumes/Data/mnt/cache
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
${CONDA_RUN} python -m pytest -n 0 -s -v --make-reports=tests_torch_mps tests/

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@@ -1,4 +1,4 @@
name: Slow tests on main
name: Run all tests
on:
push:
@@ -10,7 +10,7 @@ env:
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
PYTEST_TIMEOUT: 1000
RUN_SLOW: yes
jobs:
@@ -61,8 +61,8 @@ jobs:
- name: Install dependencies
run: |
python -m pip install -e .[quality,test]
python -m pip install -U git+https://github.com/huggingface/transformers
python -m pip install git+https://github.com/huggingface/accelerate
python -m pip install -U git+https://github.com/huggingface/transformers
- name: Environment
run: |

5
.gitignore vendored
View File

@@ -165,7 +165,4 @@ tags
# DS_Store (MacOS)
.DS_Store
# RL pipelines may produce mp4 outputs
*.mp4
# dependencies
/transformers
*.mp4

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@@ -1,40 +0,0 @@
cff-version: 1.2.0
title: 'Diffusers: State-of-the-art diffusion models'
message: >-
If you use this software, please cite it using the
metadata from this file.
type: software
authors:
- given-names: Patrick
family-names: von Platen
- given-names: Suraj
family-names: Patil
- given-names: Anton
family-names: Lozhkov
- given-names: Pedro
family-names: Cuenca
- given-names: Nathan
family-names: Lambert
- given-names: Kashif
family-names: Rasul
- given-names: Mishig
family-names: Davaadorj
- given-names: Thomas
family-names: Wolf
repository-code: 'https://github.com/huggingface/diffusers'
abstract: >-
Diffusers provides pretrained diffusion models across
multiple modalities, such as vision and audio, and serves
as a modular toolbox for inference and training of
diffusion models.
keywords:
- deep-learning
- pytorch
- image-generation
- diffusion
- text2image
- image2image
- score-based-generative-modeling
- stable-diffusion
license: Apache-2.0
version: 0.12.1

View File

@@ -45,14 +45,12 @@ quality:
isort --check-only $(check_dirs)
flake8 $(check_dirs)
doc-builder style src/diffusers docs/source --max_len 119 --check_only --path_to_docs docs/source
python utils/check_doc_toc.py
# Format source code automatically and check is there are any problems left that need manual fixing
extra_style_checks:
python utils/custom_init_isort.py
doc-builder style src/diffusers docs/source --max_len 119 --path_to_docs docs/source
python utils/check_doc_toc.py --fix_and_overwrite
# this target runs checks on all files and potentially modifies some of them

157
README.md
View File

@@ -1,6 +1,6 @@
<p align="center">
<br>
<img src="./docs/source/en/imgs/diffusers_library.jpg" width="400"/>
<img src="https://github.com/huggingface/diffusers/raw/main/docs/source/imgs/diffusers_library.jpg" width="400"/>
<br>
<p>
<p align="center">
@@ -29,13 +29,13 @@ More precisely, 🤗 Diffusers offers:
### For PyTorch
**With `pip`** (official package)
**With `pip`**
```bash
pip install --upgrade diffusers[torch]
```
**With `conda`** (maintained by the community)
**With `conda`**
```sh
conda install -c conda-forge diffusers
@@ -79,13 +79,19 @@ In order to get started, we recommend taking a look at two notebooks:
Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [LAION](https://laion.ai/) and [RunwayML](https://runwayml.com/). It's trained on 512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/) database. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and runs on a GPU with at least 4GB VRAM.
See the [model card](https://huggingface.co/CompVis/stable-diffusion) for more information.
You need to accept the model license before downloading or using the Stable Diffusion weights. Please, visit the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5), read the license carefully and tick the checkbox if you agree. You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section](https://huggingface.co/docs/hub/security-tokens) of the documentation.
### Text-to-Image generation with Stable Diffusion
First let's install
```bash
pip install --upgrade diffusers transformers accelerate
pip install --upgrade diffusers transformers scipy
```
Run this command to log in with your HF Hub token if you haven't before (you can skip this step if you prefer to run the model locally, follow [this](#running-the-model-locally) instead)
```bash
huggingface-cli login
```
We recommend using the model in [half-precision (`fp16`)](https://pytorch.org/blog/accelerating-training-on-nvidia-gpus-with-pytorch-automatic-mixed-precision/) as it gives almost always the same results as full
@@ -95,7 +101,7 @@ precision while being roughly twice as fast and requiring half the amount of GPU
import torch
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, revision="fp16")
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
@@ -103,16 +109,17 @@ image = pipe(prompt).images[0]
```
#### Running the model locally
You can also simply download the model folder and pass the path to the local folder to the `StableDiffusionPipeline`.
If you don't want to login to Hugging Face, you can also simply download the model folder
(after having [accepted the license](https://huggingface.co/runwayml/stable-diffusion-v1-5)) and pass
the path to the local folder to the `StableDiffusionPipeline`.
```
git lfs install
git clone https://huggingface.co/runwayml/stable-diffusion-v1-5
```
Assuming the folder is stored locally under `./stable-diffusion-v1-5`, you can run stable diffusion
as follows:
Assuming the folder is stored locally under `./stable-diffusion-v1-5`, you can also run stable diffusion
without requiring an authentication token:
```python
pipe = StableDiffusionPipeline.from_pretrained("./stable-diffusion-v1-5")
@@ -127,7 +134,11 @@ to using `fp16`.
The following snippet should result in less than 4GB VRAM.
```python
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
@@ -153,6 +164,7 @@ If you want to run Stable Diffusion on CPU or you want to have maximum precision
please run the model in the default *full-precision* setting:
```python
# make sure you're logged in with `huggingface-cli login`
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
@@ -235,102 +247,6 @@ images = pipeline(prompt_ids, params, prng_seed, num_inference_steps, jit=True).
images = pipeline.numpy_to_pil(np.asarray(images.reshape((num_samples,) + images.shape[-3:])))
```
Diffusers also has a Image-to-Image generation pipeline with Flax/Jax
```python
import jax
import numpy as np
import jax.numpy as jnp
from flax.jax_utils import replicate
from flax.training.common_utils import shard
import requests
from io import BytesIO
from PIL import Image
from diffusers import FlaxStableDiffusionImg2ImgPipeline
def create_key(seed=0):
return jax.random.PRNGKey(seed)
rng = create_key(0)
url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
response = requests.get(url)
init_img = Image.open(BytesIO(response.content)).convert("RGB")
init_img = init_img.resize((768, 512))
prompts = "A fantasy landscape, trending on artstation"
pipeline, params = FlaxStableDiffusionImg2ImgPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4", revision="flax",
dtype=jnp.bfloat16,
)
num_samples = jax.device_count()
rng = jax.random.split(rng, jax.device_count())
prompt_ids, processed_image = pipeline.prepare_inputs(prompt=[prompts]*num_samples, image = [init_img]*num_samples)
p_params = replicate(params)
prompt_ids = shard(prompt_ids)
processed_image = shard(processed_image)
output = pipeline(
prompt_ids=prompt_ids,
image=processed_image,
params=p_params,
prng_seed=rng,
strength=0.75,
num_inference_steps=50,
jit=True,
height=512,
width=768).images
output_images = pipeline.numpy_to_pil(np.asarray(output.reshape((num_samples,) + output.shape[-3:])))
```
Diffusers also has a Text-guided inpainting pipeline with Flax/Jax
```python
import jax
import numpy as np
from flax.jax_utils import replicate
from flax.training.common_utils import shard
import PIL
import requests
from io import BytesIO
from diffusers import FlaxStableDiffusionInpaintPipeline
def download_image(url):
response = requests.get(url)
return PIL.Image.open(BytesIO(response.content)).convert("RGB")
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
pipeline, params = FlaxStableDiffusionInpaintPipeline.from_pretrained("xvjiarui/stable-diffusion-2-inpainting")
prompt = "Face of a yellow cat, high resolution, sitting on a park bench"
prng_seed = jax.random.PRNGKey(0)
num_inference_steps = 50
num_samples = jax.device_count()
prompt = num_samples * [prompt]
init_image = num_samples * [init_image]
mask_image = num_samples * [mask_image]
prompt_ids, processed_masked_images, processed_masks = pipeline.prepare_inputs(prompt, init_image, mask_image)
# shard inputs and rng
params = replicate(params)
prng_seed = jax.random.split(prng_seed, jax.device_count())
prompt_ids = shard(prompt_ids)
processed_masked_images = shard(processed_masked_images)
processed_masks = shard(processed_masks)
images = pipeline(prompt_ids, processed_masks, processed_masked_images, params, prng_seed, num_inference_steps, jit=True).images
images = pipeline.numpy_to_pil(np.asarray(images.reshape((num_samples,) + images.shape[-3:])))
```
### Image-to-Image text-guided generation with Stable Diffusion
The `StableDiffusionImg2ImgPipeline` lets you pass a text prompt and an initial image to condition the generation of new images.
@@ -346,8 +262,11 @@ from diffusers import StableDiffusionImg2ImgPipeline
# load the pipeline
device = "cuda"
model_id_or_path = "runwayml/stable-diffusion-v1-5"
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(model_id_or_path, torch_dtype=torch.float16)
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
model_id_or_path,
revision="fp16",
torch_dtype=torch.float16,
)
# or download via git clone https://huggingface.co/runwayml/stable-diffusion-v1-5
# and pass `model_id_or_path="./stable-diffusion-v1-5"`.
pipe = pipe.to(device)
@@ -361,7 +280,7 @@ init_image = init_image.resize((768, 512))
prompt = "A fantasy landscape, trending on artstation"
images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
images[0].save("fantasy_landscape.png")
```
@@ -369,7 +288,10 @@ You can also run this example on colab [![Open In Colab](https://colab.research.
### In-painting using Stable Diffusion
The `StableDiffusionInpaintPipeline` lets you edit specific parts of an image by providing a mask and a text prompt.
The `StableDiffusionInpaintPipeline` lets you edit specific parts of an image by providing a mask and a text prompt. It uses a model optimized for this particular task, whose license you need to accept before use.
Please, visit the [model card](https://huggingface.co/runwayml/stable-diffusion-inpainting), read the license carefully and tick the checkbox if you agree. Note that this is an additional license, you need to accept it even if you accepted the text-to-image Stable Diffusion license in the past. You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section](https://huggingface.co/docs/hub/security-tokens) of the documentation.
```python
import PIL
@@ -389,7 +311,11 @@ mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data
init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
pipe = StableDiffusionInpaintPipeline.from_pretrained("runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16)
pipe = StableDiffusionInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
prompt = "Face of a yellow cat, high resolution, sitting on a park bench"
@@ -398,8 +324,11 @@ image = pipe(prompt=prompt, image=init_image, mask_image=mask_image).images[0]
### Tweak prompts reusing seeds and latents
You can generate your own latents to reproduce results, or tweak your prompt on a specific result you liked.
Please have a look at [Reusing seeds for deterministic generation](https://huggingface.co/docs/diffusers/main/en/using-diffusers/reusing_seeds).
You can generate your own latents to reproduce results, or tweak your prompt on a specific result you liked. [This notebook](https://github.com/pcuenca/diffusers-examples/blob/main/notebooks/stable-diffusion-seeds.ipynb) shows how to do it step by step. You can also run it in Google Colab [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/pcuenca/diffusers-examples/blob/main/notebooks/stable-diffusion-seeds.ipynb).
For more details, check out [the Stable Diffusion notebook](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb) [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb)
and have a look into the [release notes](https://github.com/huggingface/diffusers/releases/tag/v0.2.0).
## Fine-Tuning Stable Diffusion

View File

@@ -11,7 +11,6 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -34,8 +33,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
modelcards \
numpy \
scipy \
tensorboard \

View File

@@ -11,7 +11,6 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -36,8 +35,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
modelcards \
numpy \
scipy \
tensorboard \

View File

@@ -11,7 +11,6 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -34,8 +33,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
modelcards \
numpy \
scipy \
tensorboard \

View File

@@ -11,7 +11,6 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -34,8 +33,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
modelcards \
numpy \
scipy \
tensorboard \

View File

@@ -11,7 +11,6 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -33,8 +32,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
modelcards \
numpy \
scipy \
tensorboard \

View File

@@ -11,7 +11,6 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -33,8 +32,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
Jinja2 \
librosa \
modelcards \
numpy \
scipy \
tensorboard \

View File

@@ -1,271 +0,0 @@
<!---
Copyright 2022- The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License");
you may not use this file except in compliance with the License.
You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software
distributed under the License is distributed on an "AS IS" BASIS,
WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
See the License for the specific language governing permissions and
limitations under the License.
-->
# Generating the documentation
To generate the documentation, you first have to build it. Several packages are necessary to build the doc,
you can install them with the following command, at the root of the code repository:
```bash
pip install -e ".[docs]"
```
Then you need to install our open source documentation builder tool:
```bash
pip install git+https://github.com/huggingface/doc-builder
```
---
**NOTE**
You only need to generate the documentation to inspect it locally (if you're planning changes and want to
check how they look before committing for instance). You don't have to commit the built documentation.
---
## Previewing the documentation
To preview the docs, first install the `watchdog` module with:
```bash
pip install watchdog
```
Then run the following command:
```bash
doc-builder preview {package_name} {path_to_docs}
```
For example:
```bash
doc-builder preview diffusers docs/source/en
```
The docs will be viewable at [http://localhost:3000](http://localhost:3000). You can also preview the docs once you have opened a PR. You will see a bot add a comment to a link where the documentation with your changes lives.
---
**NOTE**
The `preview` command only works with existing doc files. When you add a completely new file, you need to update `_toctree.yml` & restart `preview` command (`ctrl-c` to stop it & call `doc-builder preview ...` again).
---
## Adding a new element to the navigation bar
Accepted files are Markdown (.md or .mdx).
Create a file with its extension and put it in the source directory. You can then link it to the toc-tree by putting
the filename without the extension in the [`_toctree.yml`](https://github.com/huggingface/diffusers/blob/main/docs/source/_toctree.yml) file.
## Renaming section headers and moving sections
It helps to keep the old links working when renaming the section header and/or moving sections from one document to another. This is because the old links are likely to be used in Issues, Forums, and Social media and it'd make for a much more superior user experience if users reading those months later could still easily navigate to the originally intended information.
Therefore, we simply keep a little map of moved sections at the end of the document where the original section was. The key is to preserve the original anchor.
So if you renamed a section from: "Section A" to "Section B", then you can add at the end of the file:
```
Sections that were moved:
[ <a href="#section-b">Section A</a><a id="section-a"></a> ]
```
and of course, if you moved it to another file, then:
```
Sections that were moved:
[ <a href="../new-file#section-b">Section A</a><a id="section-a"></a> ]
```
Use the relative style to link to the new file so that the versioned docs continue to work.
For an example of a rich moved section set please see the very end of [the transformers Trainer doc](https://github.com/huggingface/transformers/blob/main/docs/source/en/main_classes/trainer.mdx).
## Writing Documentation - Specification
The `huggingface/diffusers` documentation follows the
[Google documentation](https://sphinxcontrib-napoleon.readthedocs.io/en/latest/example_google.html) style for docstrings,
although we can write them directly in Markdown.
### Adding a new tutorial
Adding a new tutorial or section is done in two steps:
- Add a new file under `docs/source`. This file can either be ReStructuredText (.rst) or Markdown (.md).
- Link that file in `docs/source/_toctree.yml` on the correct toc-tree.
Make sure to put your new file under the proper section. It's unlikely to go in the first section (*Get Started*), so
depending on the intended targets (beginners, more advanced users, or researchers) it should go in sections two, three, or four.
### Adding a new pipeline/scheduler
When adding a new pipeline:
- create a file `xxx.mdx` under `docs/source/api/pipelines` (don't hesitate to copy an existing file as template).
- Link that file in (*Diffusers Summary*) section in `docs/source/api/pipelines/overview.mdx`, along with the link to the paper, and a colab notebook (if available).
- Write a short overview of the diffusion model:
- Overview with paper & authors
- Paper abstract
- Tips and tricks and how to use it best
- Possible an end-to-end example of how to use it
- Add all the pipeline classes that should be linked in the diffusion model. These classes should be added using our Markdown syntax. By default as follows:
```
## XXXPipeline
[[autodoc]] XXXPipeline
- all
- __call__
```
This will include every public method of the pipeline that is documented, as well as the `__call__` method that is not documented by default. If you just want to add additional methods that are not documented, you can put the list of all methods to add in a list that contains `all`.
```
[[autodoc]] XXXPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
```
You can follow the same process to create a new scheduler under the `docs/source/api/schedulers` folder
### Writing source documentation
Values that should be put in `code` should either be surrounded by backticks: \`like so\`. Note that argument names
and objects like True, None, or any strings should usually be put in `code`.
When mentioning a class, function, or method, it is recommended to use our syntax for internal links so that our tool
adds a link to its documentation with this syntax: \[\`XXXClass\`\] or \[\`function\`\]. This requires the class or
function to be in the main package.
If you want to create a link to some internal class or function, you need to
provide its path. For instance: \[\`pipelines.ImagePipelineOutput\`\]. This will be converted into a link with
`pipelines.ImagePipelineOutput` in the description. To get rid of the path and only keep the name of the object you are
linking to in the description, add a ~: \[\`~pipelines.ImagePipelineOutput\`\] will generate a link with `ImagePipelineOutput` in the description.
The same works for methods so you can either use \[\`XXXClass.method\`\] or \[~\`XXXClass.method\`\].
#### Defining arguments in a method
Arguments should be defined with the `Args:` (or `Arguments:` or `Parameters:`) prefix, followed by a line return and
an indentation. The argument should be followed by its type, with its shape if it is a tensor, a colon, and its
description:
```
Args:
n_layers (`int`): The number of layers of the model.
```
If the description is too long to fit in one line, another indentation is necessary before writing the description
after the argument.
Here's an example showcasing everything so far:
```
Args:
input_ids (`torch.LongTensor` of shape `(batch_size, sequence_length)`):
Indices of input sequence tokens in the vocabulary.
Indices can be obtained using [`AlbertTokenizer`]. See [`~PreTrainedTokenizer.encode`] and
[`~PreTrainedTokenizer.__call__`] for details.
[What are input IDs?](../glossary#input-ids)
```
For optional arguments or arguments with defaults we follow the following syntax: imagine we have a function with the
following signature:
```
def my_function(x: str = None, a: float = 1):
```
then its documentation should look like this:
```
Args:
x (`str`, *optional*):
This argument controls ...
a (`float`, *optional*, defaults to 1):
This argument is used to ...
```
Note that we always omit the "defaults to \`None\`" when None is the default for any argument. Also note that even
if the first line describing your argument type and its default gets long, you can't break it on several lines. You can
however write as many lines as you want in the indented description (see the example above with `input_ids`).
#### Writing a multi-line code block
Multi-line code blocks can be useful for displaying examples. They are done between two lines of three backticks as usual in Markdown:
````
```
# first line of code
# second line
# etc
```
````
#### Writing a return block
The return block should be introduced with the `Returns:` prefix, followed by a line return and an indentation.
The first line should be the type of the return, followed by a line return. No need to indent further for the elements
building the return.
Here's an example of a single value return:
```
Returns:
`List[int]`: A list of integers in the range [0, 1] --- 1 for a special token, 0 for a sequence token.
```
Here's an example of a tuple return, comprising several objects:
```
Returns:
`tuple(torch.FloatTensor)` comprising various elements depending on the configuration ([`BertConfig`]) and inputs:
- ** loss** (*optional*, returned when `masked_lm_labels` is provided) `torch.FloatTensor` of shape `(1,)` --
Total loss is the sum of the masked language modeling loss and the next sequence prediction (classification) loss.
- **prediction_scores** (`torch.FloatTensor` of shape `(batch_size, sequence_length, config.vocab_size)`) --
Prediction scores of the language modeling head (scores for each vocabulary token before SoftMax).
```
#### Adding an image
Due to the rapidly growing repository, it is important to make sure that no files that would significantly weigh down the repository are added. This includes images, videos, and other non-text files. We prefer to leverage a hf.co hosted `dataset` like
the ones hosted on [`hf-internal-testing`](https://huggingface.co/hf-internal-testing) in which to place these files and reference
them by URL. We recommend putting them in the following dataset: [huggingface/documentation-images](https://huggingface.co/datasets/huggingface/documentation-images).
If an external contribution, feel free to add the images to your PR and ask a Hugging Face member to migrate your images
to this dataset.
## Styling the docstring
We have an automatic script running with the `make style` command that will make sure that:
- the docstrings fully take advantage of the line width
- all code examples are formatted using black, like the code of the Transformers library
This script may have some weird failures if you made a syntax mistake or if you uncover a bug. Therefore, it's
recommended to commit your changes before running `make style`, so you can revert the changes done by that script
easily.

View File

@@ -1,57 +0,0 @@
### Translating the Diffusers documentation into your language
As part of our mission to democratize machine learning, we'd love to make the Diffusers library available in many more languages! Follow the steps below if you want to help translate the documentation into your language 🙏.
**🗞️ Open an issue**
To get started, navigate to the [Issues](https://github.com/huggingface/diffusers/issues) page of this repo and check if anyone else has opened an issue for your language. If not, open a new issue by selecting the "Translation template" from the "New issue" button.
Once an issue exists, post a comment to indicate which chapters you'd like to work on, and we'll add your name to the list.
**🍴 Fork the repository**
First, you'll need to [fork the Diffusers repo](https://docs.github.com/en/get-started/quickstart/fork-a-repo). You can do this by clicking on the **Fork** button on the top-right corner of this repo's page.
Once you've forked the repo, you'll want to get the files on your local machine for editing. You can do that by cloning the fork with Git as follows:
```bash
git clone https://github.com/YOUR-USERNAME/diffusers.git
```
**📋 Copy-paste the English version with a new language code**
The documentation files are in one leading directory:
- [`docs/source`](https://github.com/huggingface/diffusers/tree/main/docs/source): All the documentation materials are organized here by language.
You'll only need to copy the files in the [`docs/source/en`](https://github.com/huggingface/diffusers/tree/main/docs/source/en) directory, so first navigate to your fork of the repo and run the following:
```bash
cd ~/path/to/diffusers/docs
cp -r source/en source/LANG-ID
```
Here, `LANG-ID` should be one of the ISO 639-1 or ISO 639-2 language codes -- see [here](https://www.loc.gov/standards/iso639-2/php/code_list.php) for a handy table.
**✍️ Start translating**
The fun part comes - translating the text!
The first thing we recommend is translating the part of the `_toctree.yml` file that corresponds to your doc chapter. This file is used to render the table of contents on the website.
> 🙋 If the `_toctree.yml` file doesn't yet exist for your language, you can create one by copy-pasting from the English version and deleting the sections unrelated to your chapter. Just make sure it exists in the `docs/source/LANG-ID/` directory!
The fields you should add are `local` (with the name of the file containing the translation; e.g. `autoclass_tutorial`), and `title` (with the title of the doc in your language; e.g. `Load pretrained instances with an AutoClass`) -- as a reference, here is the `_toctree.yml` for [English](https://github.com/huggingface/diffusers/blob/main/docs/source/en/_toctree.yml):
```yaml
- sections:
- local: pipeline_tutorial # Do not change this! Use the same name for your .md file
title: Pipelines for inference # Translate this!
...
title: Tutorials # Translate this!
```
Once you have translated the `_toctree.yml` file, you can start translating the [MDX](https://mdxjs.com/) files associated with your docs chapter.
> 🙋 If you'd like others to help you with the translation, you should [open an issue](https://github.com/huggingface/diffusers/issues) and tag @patrickvonplaten.

126
docs/source/_toctree.yml Normal file
View File

@@ -0,0 +1,126 @@
- sections:
- local: index
title: "🧨 Diffusers"
- local: quicktour
title: "Quicktour"
- local: installation
title: "Installation"
title: "Get started"
- sections:
- sections:
- local: using-diffusers/loading
title: "Loading Pipelines, Models, and Schedulers"
- local: using-diffusers/schedulers
title: "Using different Schedulers"
- local: using-diffusers/configuration
title: "Configuring Pipelines, Models, and Schedulers"
- local: using-diffusers/custom_pipeline_overview
title: "Loading and Adding Custom Pipelines"
title: "Loading & Hub"
- sections:
- local: using-diffusers/unconditional_image_generation
title: "Unconditional Image Generation"
- local: using-diffusers/conditional_image_generation
title: "Text-to-Image Generation"
- local: using-diffusers/img2img
title: "Text-Guided Image-to-Image"
- local: using-diffusers/inpaint
title: "Text-Guided Image-Inpainting"
- local: using-diffusers/custom_pipeline_examples
title: "Community Pipelines"
- local: using-diffusers/contribute_pipeline
title: "How to contribute a Pipeline"
title: "Pipelines for Inference"
- sections:
- local: using-diffusers/rl
title: "Reinforcement Learning"
- local: using-diffusers/audio
title: "Audio"
- local: using-diffusers/other-modalities
title: "Other Modalities"
title: "Taking Diffusers Beyond Images"
title: "Using Diffusers"
- sections:
- local: optimization/fp16
title: "Memory and Speed"
- local: optimization/onnx
title: "ONNX"
- local: optimization/open_vino
title: "OpenVINO"
- local: optimization/mps
title: "MPS"
title: "Optimization/Special Hardware"
- sections:
- local: training/overview
title: "Overview"
- local: training/unconditional_training
title: "Unconditional Image Generation"
- local: training/text_inversion
title: "Textual Inversion"
- local: training/dreambooth
title: "Dreambooth"
- local: training/text2image
title: "Text-to-image fine-tuning"
title: "Training"
- sections:
- local: conceptual/stable_diffusion
title: "Stable Diffusion"
- local: conceptual/philosophy
title: "Philosophy"
- local: conceptual/contribution
title: "How to contribute?"
title: "Conceptual Guides"
- sections:
- sections:
- local: api/models
title: "Models"
- local: api/schedulers
title: "Schedulers"
- local: api/diffusion_pipeline
title: "Diffusion Pipeline"
- local: api/logging
title: "Logging"
- local: api/configuration
title: "Configuration"
- local: api/outputs
title: "Outputs"
title: "Main Classes"
- sections:
- local: api/pipelines/overview
title: "Overview"
- local: api/pipelines/alt_diffusion
title: "AltDiffusion"
- local: api/pipelines/cycle_diffusion
title: "Cycle Diffusion"
- local: api/pipelines/ddim
title: "DDIM"
- local: api/pipelines/ddpm
title: "DDPM"
- local: api/pipelines/latent_diffusion
title: "Latent Diffusion"
- local: api/pipelines/latent_diffusion_uncond
title: "Unconditional Latent Diffusion"
- local: api/pipelines/pndm
title: "PNDM"
- local: api/pipelines/score_sde_ve
title: "Score SDE VE"
- local: api/pipelines/stable_diffusion
title: "Stable Diffusion"
- local: api/pipelines/stable_diffusion_safe
title: "Safe Stable Diffusion"
- local: api/pipelines/stochastic_karras_ve
title: "Stochastic Karras VE"
- local: api/pipelines/dance_diffusion
title: "Dance Diffusion"
- local: api/pipelines/versatile_diffusion
title: "Versatile Diffusion"
- local: api/pipelines/vq_diffusion
title: "VQ Diffusion"
- local: api/pipelines/repaint
title: "RePaint"
title: "Pipelines"
- sections:
- local: api/experimental/rl
title: "RL Planning"
title: "Experimental Features"
title: "API"

View File

@@ -30,18 +30,13 @@ Any pipeline object can be saved locally with [`~DiffusionPipeline.save_pretrain
## DiffusionPipeline
[[autodoc]] DiffusionPipeline
- all
- __call__
- device
- from_pretrained
- save_pretrained
- to
- device
- components
## ImagePipelineOutput
By default diffusion pipelines return an object of class
[[autodoc]] pipelines.ImagePipelineOutput
## AudioPipelineOutput
By default diffusion pipelines return an object of class
[[autodoc]] pipelines.AudioPipelineOutput
[[autodoc]] pipeline_utils.ImagePipelineOutput

View File

@@ -1,4 +1,4 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
<!--Copyright 2020 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at

View File

@@ -41,13 +41,13 @@ The models are built on the base class ['ModelMixin'] that is a `torch.nn.module
[[autodoc]] models.vae.DecoderOutput
## VQEncoderOutput
[[autodoc]] models.vq_model.VQEncoderOutput
[[autodoc]] models.vae.VQEncoderOutput
## VQModel
[[autodoc]] VQModel
## AutoencoderKLOutput
[[autodoc]] models.autoencoder_kl.AutoencoderKLOutput
[[autodoc]] models.vae.AutoencoderKLOutput
## AutoencoderKL
[[autodoc]] AutoencoderKL
@@ -56,13 +56,7 @@ The models are built on the base class ['ModelMixin'] that is a `torch.nn.module
[[autodoc]] Transformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.transformer_2d.Transformer2DModelOutput
## PriorTransformer
[[autodoc]] models.prior_transformer.PriorTransformer
## PriorTransformerOutput
[[autodoc]] models.prior_transformer.PriorTransformerOutput
[[autodoc]] models.attention.Transformer2DModelOutput
## FlaxModelMixin
[[autodoc]] FlaxModelMixin

View File

@@ -25,7 +25,7 @@ pipeline = DDIMPipeline.from_pretrained("google/ddpm-cifar10-32")
outputs = pipeline()
```
The `outputs` object is a [`~pipelines.ImagePipelineOutput`], as we can see in the
The `outputs` object is a [`~pipeline_utils.ImagePipelineOutput`], as we can see in the
documentation of that class below, it means it has an image attribute.
You can access each attribute as you would usually do, and if that attribute has not been returned by the model, you will get `None`:

View File

@@ -28,7 +28,7 @@ The abstract of the paper is the following:
## Tips
- AltDiffusion is conceptually exaclty the same as [Stable Diffusion](./api/pipelines/stable_diffusion/overview).
- AltDiffusion is conceptually exaclty the same as [Stable Diffusion](./api/pipelines/stable_diffusion).
- *Run AltDiffusion*
@@ -51,7 +51,7 @@ To use a different scheduler, you can either change it via the [`ConfigMixin.fro
```
- *How to convert all use cases with multiple or single pipeline*
- *How to conver all use cases with multiple or single pipeline*
If you want to use all possible use cases in a single `DiffusionPipeline` we recommend using the `components` functionality to instantiate all components in the most memory-efficient way:
@@ -69,15 +69,15 @@ If you want to use all possible use cases in a single `DiffusionPipeline` we rec
## AltDiffusionPipelineOutput
[[autodoc]] pipelines.alt_diffusion.AltDiffusionPipelineOutput
- all
- __call__
## AltDiffusionPipeline
[[autodoc]] AltDiffusionPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
## AltDiffusionImg2ImgPipeline
[[autodoc]] AltDiffusionImg2ImgPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing

View File

@@ -57,7 +57,7 @@ prompt = "An astronaut riding an elephant"
image = pipe(
prompt=prompt,
source_prompt=source_prompt,
image=init_image,
init_image=init_image,
num_inference_steps=100,
eta=0.1,
strength=0.8,
@@ -83,7 +83,7 @@ torch.manual_seed(0)
image = pipe(
prompt=prompt,
source_prompt=source_prompt,
image=init_image,
init_image=init_image,
num_inference_steps=100,
eta=0.1,
strength=0.85,
@@ -96,5 +96,4 @@ image.save("black_to_blue.png")
## CycleDiffusionPipeline
[[autodoc]] CycleDiffusionPipeline
- all
- __call__

View File

@@ -30,5 +30,4 @@ The original codebase of this implementation can be found [here](https://github.
## DanceDiffusionPipeline
[[autodoc]] DanceDiffusionPipeline
- all
- __call__
- __call__

View File

@@ -32,5 +32,4 @@ For questions, feel free to contact the author on [tsong.me](https://tsong.me/).
## DDIMPipeline
[[autodoc]] DDIMPipeline
- all
- __call__
- __call__

View File

@@ -33,5 +33,4 @@ The original codebase of this paper can be found [here](https://github.com/hojon
# DDPMPipeline
[[autodoc]] DDPMPipeline
- all
- __call__
- __call__

View File

@@ -40,10 +40,8 @@ The original codebase can be found [here](https://github.com/CompVis/latent-diff
## LDMTextToImagePipeline
[[autodoc]] LDMTextToImagePipeline
- all
- __call__
- __call__
## LDMSuperResolutionPipeline
[[autodoc]] LDMSuperResolutionPipeline
- all
- __call__
- __call__

View File

@@ -38,5 +38,4 @@ The original codebase can be found [here](https://github.com/CompVis/latent-diff
## LDMPipeline
[[autodoc]] LDMPipeline
- all
- __call__
- __call__

View File

@@ -44,32 +44,26 @@ available a colab notebook to directly try them out.
| Pipeline | Paper | Tasks | Colab
|---|---|:---:|:---:|
| [alt_diffusion](./alt_diffusion) | [**AltDiffusion**](https://arxiv.org/abs/2211.06679) | Image-to-Image Text-Guided Generation | -
| [audio_diffusion](./audio_diffusion) | [**Audio Diffusion**](https://github.com/teticio/audio_diffusion.git) | Unconditional Audio Generation |
| [cycle_diffusion](./cycle_diffusion) | [**Cycle Diffusion**](https://arxiv.org/abs/2210.05559) | Image-to-Image Text-Guided Generation |
| [dance_diffusion](./dance_diffusion) | [**Dance Diffusion**](https://github.com/williamberman/diffusers.git) | Unconditional Audio Generation |
| [ddpm](./ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) | Unconditional Image Generation |
| [ddim](./ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) | Unconditional Image Generation |
| [latent_diffusion](./latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation |
| [latent_diffusion](./latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image |
| [latent_diffusion_uncond](./latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation |
| [paint_by_example](./paint_by_example) | [**Paint by Example: Exemplar-based Image Editing with Diffusion Models**](https://arxiv.org/abs/2211.13227) | Image-Guided Image Inpainting |
| [pndm](./pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation |
| [score_sde_ve](./score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [score_sde_vp](./score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [stable_diffusion](./stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
| [stable_diffusion](./stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
| [stable_diffusion](./stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-to-Image Generation |
| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting |
| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
| [stable_diffusion_safe](./stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
| [stochastic_karras_ve](./stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
| [unclip](./unclip) | [Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) | Text-to-Image Generation |
| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation |
| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |
| [vq_diffusion](./vq_diffusion) | [Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation |
| [alt_diffusion](./api/pipelines/alt_diffusion) | [**AltDiffusion**](https://arxiv.org/abs/2211.06679) | Image-to-Image Text-Guided Generation | -
| [cycle_diffusion](./api/pipelines/cycle_diffusion) | [**Cycle Diffusion**](https://arxiv.org/abs/2210.05559) | Image-to-Image Text-Guided Generation |
| [dance_diffusion](./api/pipelines/dance_diffusion) | [**Dance Diffusion**](https://github.com/williamberman/diffusers.git) | Unconditional Audio Generation |
| [ddpm](./api/pipelines/ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) | Unconditional Image Generation |
| [ddim](./api/pipelines/ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) | Unconditional Image Generation |
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation |
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image |
| [latent_diffusion_uncond](./api/pipelines/latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation |
| [pndm](./api/pipelines/pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation |
| [score_sde_ve](./api/pipelines/score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [score_sde_vp](./api/pipelines/score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
| [stable_diffusion_safe](./api/pipelines/stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
| [stochastic_karras_ve](./api/pipelines/stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |
| [vq_diffusion](./api/pipelines/vq_diffusion) | [Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation |
**Note**: Pipelines are simple examples of how to play around with the diffusion systems as described in the corresponding papers.
@@ -139,9 +133,9 @@ from diffusers import StableDiffusionImg2ImgPipeline
# load the pipeline
device = "cuda"
pipe = StableDiffusionImg2ImgPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16).to(
device
)
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", revision="fp16", torch_dtype=torch.float16
).to(device)
# let's download an initial image
url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
@@ -152,7 +146,7 @@ init_image = init_image.resize((768, 512))
prompt = "A fantasy landscape, trending on artstation"
images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
images[0].save("fantasy_landscape.png")
```
@@ -189,6 +183,7 @@ mask_image = download_image(mask_url).resize((512, 512))
pipe = StableDiffusionInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")

View File

@@ -30,6 +30,6 @@ The original codebase can be found [here](https://github.com/luping-liu/PNDM).
## PNDMPipeline
[[autodoc]] PNDMPipeline
- all
- __call__
[[autodoc]] pipelines.pndm.pipeline_pndm.PNDMPipeline
- __call__

View File

@@ -72,6 +72,6 @@ inpainted_image = output.images[0]
```
## RePaintPipeline
[[autodoc]] RePaintPipeline
- all
- __call__
[[autodoc]] pipelines.repaint.pipeline_repaint.RePaintPipeline
- __call__

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@@ -32,5 +32,5 @@ This pipeline implements the Variance Expanding (VE) variant of the method.
## ScoreSdeVePipeline
[[autodoc]] ScoreSdeVePipeline
- all
- __call__
- __call__

View File

@@ -25,16 +25,9 @@ For more details about how Stable Diffusion works and how it differs from the ba
| Pipeline | Tasks | Colab | Demo
|---|---|:---:|:---:|
| [StableDiffusionPipeline](./text2img) | *Text-to-Image Generation* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb) | [🤗 Stable Diffusion](https://huggingface.co/spaces/stabilityai/stable-diffusion)
| [StableDiffusionImg2ImgPipeline](./img2img) | *Image-to-Image Text-Guided Generation* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb) | [🤗 Diffuse the Rest](https://huggingface.co/spaces/huggingface/diffuse-the-rest)
| [StableDiffusionInpaintPipeline](./inpaint) | **Experimental** *Text-Guided Image Inpainting* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb) | Coming soon
| [StableDiffusionDepth2ImgPipeline](./depth2img) | **Experimental** *Depth-to-Image Text-Guided Generation * | | Coming soon
| [StableDiffusionImageVariationPipeline](./image_variation) | **Experimental** *Image Variation Generation * | | [🤗 Stable Diffusion Image Variations](https://huggingface.co/spaces/lambdalabs/stable-diffusion-image-variations)
| [StableDiffusionUpscalePipeline](./upscale) | **Experimental** *Text-Guided Image Super-Resolution * | | Coming soon
| [StableDiffusionLatentUpscalePipeline](./latent_upscale) | **Experimental** *Text-Guided Image Super-Resolution * | | Coming soon
| [StableDiffusionInstructPix2PixPipeline](./pix2pix) | **Experimental** *Text-Based Image Editing * | | [InstructPix2Pix: Learning to Follow Image Editing Instructions](https://huggingface.co/spaces/timbrooks/instruct-pix2pix)
| [pipeline_stable_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py) | *Text-to-Image Generation* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb) | [🤗 Stable Diffusion](https://huggingface.co/spaces/stabilityai/stable-diffusion)
| [pipeline_stable_diffusion_img2img.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_img2img.py) | *Image-to-Image Text-Guided Generation* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb) | [🤗 Diffuse the Rest](https://huggingface.co/spaces/huggingface/diffuse-the-rest)
| [pipeline_stable_diffusion_inpaint.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_inpaint.py) | **Experimental** *Text-Guided Image Inpainting* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb) | Coming soon
## Tips
@@ -55,7 +48,7 @@ To use a different scheduler, you can either change it via the [`ConfigMixin.fro
```
### How to convert all use cases with multiple or single pipeline
### How to conver all use cases with multiple or single pipeline
If you want to use all possible use cases in a single `DiffusionPipeline` you can either:
- Make use of the [Stable Diffusion Mega Pipeline](https://github.com/huggingface/diffusers/tree/main/examples/community#stable-diffusion-mega) or
@@ -77,3 +70,28 @@ If you want to use all possible use cases in a single `DiffusionPipeline` you ca
## StableDiffusionPipelineOutput
[[autodoc]] pipelines.stable_diffusion.StableDiffusionPipelineOutput
## StableDiffusionPipeline
[[autodoc]] StableDiffusionPipeline
- __call__
- enable_attention_slicing
- disable_attention_slicing
## StableDiffusionImg2ImgPipeline
[[autodoc]] StableDiffusionImg2ImgPipeline
- __call__
- enable_attention_slicing
- disable_attention_slicing
## StableDiffusionInpaintPipeline
[[autodoc]] StableDiffusionInpaintPipeline
- __call__
- enable_attention_slicing
- disable_attention_slicing
## StableDiffusionImageVariationPipeline
[[autodoc]] StableDiffusionImageVariationPipeline
- __call__
- enable_attention_slicing
- disable_attention_slicing

View File

@@ -28,7 +28,7 @@ The abstract of the paper is the following:
## Tips
- Safe Stable Diffusion may also be used with weights of [Stable Diffusion](./api/pipelines/stable_diffusion/text2img).
- Safe Stable Diffusion may also be used with weights of [Stable Diffusion](./api/pipelines/stable_diffusion).
### Run Safe Stable Diffusion
@@ -81,10 +81,10 @@ To use a different scheduler, you can either change it via the [`ConfigMixin.fro
## StableDiffusionSafePipelineOutput
[[autodoc]] pipelines.stable_diffusion_safe.StableDiffusionSafePipelineOutput
- all
- __call__
## StableDiffusionPipelineSafe
[[autodoc]] StableDiffusionPipelineSafe
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing

View File

@@ -32,5 +32,4 @@ This pipeline implements the Stochastic sampling tailored to the Variance-Expand
## KarrasVePipeline
[[autodoc]] KarrasVePipeline
- all
- __call__
- __call__

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@@ -20,7 +20,7 @@ The abstract of the paper is the following:
## Tips
- VersatileDiffusion is conceptually very similar as [Stable Diffusion](./api/pipelines/stable_diffusion/overview), but instead of providing just a image data stream conditioned on text, VersatileDiffusion provides both a image and text data stream and can be conditioned on both text and image.
- VersatileDiffusion is conceptually very similar as [Stable Diffusion](./api/pipelines/stable_diffusion), but instead of providing just a image data stream conditioned on text, VersatileDiffusion provides both a image and text data stream and can be conditioned on both text and image.
### *Run VersatileDiffusion*
@@ -56,15 +56,18 @@ To use a different scheduler, you can either change it via the [`ConfigMixin.fro
## VersatileDiffusionTextToImagePipeline
[[autodoc]] VersatileDiffusionTextToImagePipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
## VersatileDiffusionImageVariationPipeline
[[autodoc]] VersatileDiffusionImageVariationPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
## VersatileDiffusionDualGuidedPipeline
[[autodoc]] VersatileDiffusionDualGuidedPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing

View File

@@ -30,6 +30,5 @@ The original codebase can be found [here](https://github.com/microsoft/VQ-Diffus
## VQDiffusionPipeline
[[autodoc]] VQDiffusionPipeline
- all
- __call__
[[autodoc]] pipelines.vq_diffusion.pipeline_vq_diffusion.VQDiffusionPipeline
- __call__

View File

@@ -0,0 +1,151 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Schedulers
Diffusers contains multiple pre-built schedule functions for the diffusion process.
## What is a scheduler?
The schedule functions, denoted *Schedulers* in the library take in the output of a trained model, a sample which the diffusion process is iterating on, and a timestep to return a denoised sample. That's why schedulers may also be called *Samplers* in other diffusion models implementations.
- Schedulers define the methodology for iteratively adding noise to an image or for updating a sample based on model outputs.
- adding noise in different manners represent the algorithmic processes to train a diffusion model by adding noise to images.
- for inference, the scheduler defines how to update a sample based on an output from a pretrained model.
- Schedulers are often defined by a *noise schedule* and an *update rule* to solve the differential equation solution.
### Discrete versus continuous schedulers
All schedulers take in a timestep to predict the updated version of the sample being diffused.
The timesteps dictate where in the diffusion process the step is, where data is generated by iterating forward in time and inference is executed by propagating backwards through timesteps.
Different algorithms use timesteps that both discrete (accepting `int` inputs), such as the [`DDPMScheduler`] or [`PNDMScheduler`], and continuous (accepting `float` inputs), such as the score-based schedulers [`ScoreSdeVeScheduler`] or [`ScoreSdeVpScheduler`].
## Designing Re-usable schedulers
The core design principle between the schedule functions is to be model, system, and framework independent.
This allows for rapid experimentation and cleaner abstractions in the code, where the model prediction is separated from the sample update.
To this end, the design of schedulers is such that:
- Schedulers can be used interchangeably between diffusion models in inference to find the preferred trade-off between speed and generation quality.
- Schedulers are currently by default in PyTorch, but are designed to be framework independent (partial Jax support currently exists).
## API
The core API for any new scheduler must follow a limited structure.
- Schedulers should provide one or more `def step(...)` functions that should be called to update the generated sample iteratively.
- Schedulers should provide a `set_timesteps(...)` method that configures the parameters of a schedule function for a specific inference task.
- Schedulers should be framework-specific.
The base class [`SchedulerMixin`] implements low level utilities used by multiple schedulers.
### SchedulerMixin
[[autodoc]] SchedulerMixin
### SchedulerOutput
The class [`SchedulerOutput`] contains the outputs from any schedulers `step(...)` call.
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
### Implemented Schedulers
#### Denoising diffusion implicit models (DDIM)
Original paper can be found here.
[[autodoc]] DDIMScheduler
#### Denoising diffusion probabilistic models (DDPM)
Original paper can be found [here](https://arxiv.org/abs/2010.02502).
[[autodoc]] DDPMScheduler
#### Multistep DPM-Solver
Original paper can be found [here](https://arxiv.org/abs/2206.00927) and the [improved version](https://arxiv.org/abs/2211.01095). The original implementation can be found [here](https://github.com/LuChengTHU/dpm-solver).
[[autodoc]] DPMSolverMultistepScheduler
#### Variance exploding, stochastic sampling from Karras et. al
Original paper can be found [here](https://arxiv.org/abs/2006.11239).
[[autodoc]] KarrasVeScheduler
#### Linear multistep scheduler for discrete beta schedules
Original implementation can be found [here](https://arxiv.org/abs/2206.00364).
[[autodoc]] LMSDiscreteScheduler
#### Pseudo numerical methods for diffusion models (PNDM)
Original implementation can be found [here](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181).
[[autodoc]] PNDMScheduler
#### variance exploding stochastic differential equation (VE-SDE) scheduler
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
[[autodoc]] ScoreSdeVeScheduler
#### improved pseudo numerical methods for diffusion models (iPNDM)
Original implementation can be found [here](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296).
[[autodoc]] IPNDMScheduler
#### variance preserving stochastic differential equation (VP-SDE) scheduler
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
<Tip warning={true}>
Score SDE-VP is under construction.
</Tip>
[[autodoc]] schedulers.scheduling_sde_vp.ScoreSdeVpScheduler
#### Euler scheduler
Euler scheduler (Algorithm 2) from the paper [Elucidating the Design Space of Diffusion-Based Generative Models](https://arxiv.org/abs/2206.00364) by Karras et al. (2022). Based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L51) implementation by Katherine Crowson.
Fast scheduler which often times generates good outputs with 20-30 steps.
[[autodoc]] EulerDiscreteScheduler
#### Euler Ancestral scheduler
Ancestral sampling with Euler method steps. Based on the original (k-diffusion)[https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72] implementation by Katherine Crowson.
Fast scheduler which often times generates good outputs with 20-30 steps.
[[autodoc]] EulerAncestralDiscreteScheduler
#### VQDiffusionScheduler
Original paper can be found [here](https://arxiv.org/abs/2111.14822)
[[autodoc]] VQDiffusionScheduler
#### RePaint scheduler
DDPM-based inpainting scheduler for unsupervised inpainting with extreme masks.
Intended for use with [`RePaintPipeline`].
Based on the paper [RePaint: Inpainting using Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2201.09865)
and the original implementation by Andreas Lugmayr et al.: https://github.com/andreas128/RePaint
[[autodoc]] RePaintScheduler

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@@ -0,0 +1,17 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Philosophy
- Readability and clarity are preferred over highly optimized code. A strong importance is put on providing readable, intuitive and elementary code design. *E.g.*, the provided [schedulers](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers) are separated from the provided [models](https://github.com/huggingface/diffusers/tree/main/src/diffusers/models) and use well-commented code that can be read alongside the original paper.
- Diffusers is **modality independent** and focuses on providing pretrained models and tools to build systems that generate **continuous outputs**, *e.g.* vision and audio. This is one of the guiding goals even if the initial pipelines are devoted to vision tasks.
- Diffusion models and schedulers are provided as concise, elementary building blocks. In contrast, diffusion pipelines are a collection of end-to-end diffusion systems that can be used out-of-the-box, should stay as close as possible to their original implementations and can include components of other libraries, such as text encoders. Examples of diffusion pipelines are [Glide](https://github.com/openai/glide-text2im), [Latent Diffusion](https://github.com/CompVis/latent-diffusion) and [Stable Diffusion](https://github.com/compvis/stable-diffusion).

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@@ -10,7 +10,6 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# 번역중
# Stable Diffusion
열심히 번역을 진행중입니다. 조금만 기다려주세요.
감사합니다!
Please visit this [very in-detail blog post](https://huggingface.co/blog/stable_diffusion) on Stable Diffusion!

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- sections:
- local: index
title: 🧨 Diffusers
- local: quicktour
title: Quicktour
- local: stable_diffusion
title: Stable Diffusion
- local: installation
title: Installation
title: Get started
- sections:
- sections:
- local: using-diffusers/loading
title: Loading Pipelines, Models, and Schedulers
- local: using-diffusers/schedulers
title: Using different Schedulers
- local: using-diffusers/configuration
title: Configuring Pipelines, Models, and Schedulers
- local: using-diffusers/custom_pipeline_overview
title: Loading and Adding Custom Pipelines
- local: using-diffusers/kerascv
title: Using KerasCV Stable Diffusion Checkpoints in Diffusers
title: Loading & Hub
- sections:
- local: using-diffusers/unconditional_image_generation
title: Unconditional Image Generation
- local: using-diffusers/conditional_image_generation
title: Text-to-Image Generation
- local: using-diffusers/img2img
title: Text-Guided Image-to-Image
- local: using-diffusers/inpaint
title: Text-Guided Image-Inpainting
- local: using-diffusers/depth2img
title: Text-Guided Depth-to-Image
- local: using-diffusers/reusing_seeds
title: Reusing seeds for deterministic generation
- local: using-diffusers/reproducibility
title: Reproducibility
- local: using-diffusers/custom_pipeline_examples
title: Community Pipelines
- local: using-diffusers/contribute_pipeline
title: How to contribute a Pipeline
- local: using-diffusers/using_safetensors
title: Using safetensors
title: Pipelines for Inference
- sections:
- local: using-diffusers/rl
title: Reinforcement Learning
- local: using-diffusers/audio
title: Audio
- local: using-diffusers/other-modalities
title: Other Modalities
title: Taking Diffusers Beyond Images
title: Using Diffusers
- sections:
- local: optimization/fp16
title: Memory and Speed
- local: optimization/xformers
title: xFormers
- local: optimization/onnx
title: ONNX
- local: optimization/open_vino
title: OpenVINO
- local: optimization/mps
title: MPS
- local: optimization/habana
title: Habana Gaudi
title: Optimization/Special Hardware
- sections:
- local: training/overview
title: Overview
- local: training/unconditional_training
title: Unconditional Image Generation
- local: training/text_inversion
title: Textual Inversion
- local: training/dreambooth
title: Dreambooth
- local: training/text2image
title: Text-to-image fine-tuning
- local: training/lora
title: LoRA Support in Diffusers
title: Training
- sections:
- local: conceptual/philosophy
title: Philosophy
- local: conceptual/contribution
title: How to contribute?
title: Conceptual Guides
- sections:
- sections:
- local: api/models
title: Models
- local: api/diffusion_pipeline
title: Diffusion Pipeline
- local: api/logging
title: Logging
- local: api/configuration
title: Configuration
- local: api/outputs
title: Outputs
- local: api/loaders
title: Loaders
title: Main Classes
- sections:
- local: api/pipelines/overview
title: Overview
- local: api/pipelines/alt_diffusion
title: AltDiffusion
- local: api/pipelines/audio_diffusion
title: Audio Diffusion
- local: api/pipelines/cycle_diffusion
title: Cycle Diffusion
- local: api/pipelines/dance_diffusion
title: Dance Diffusion
- local: api/pipelines/ddim
title: DDIM
- local: api/pipelines/ddpm
title: DDPM
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/latent_diffusion
title: Latent Diffusion
- local: api/pipelines/paint_by_example
title: PaintByExample
- local: api/pipelines/pndm
title: PNDM
- local: api/pipelines/repaint
title: RePaint
- local: api/pipelines/stable_diffusion_safe
title: Safe Stable Diffusion
- local: api/pipelines/score_sde_ve
title: Score SDE VE
- sections:
- local: api/pipelines/stable_diffusion/overview
title: Overview
- local: api/pipelines/stable_diffusion/text2img
title: Text-to-Image
- local: api/pipelines/stable_diffusion/img2img
title: Image-to-Image
- local: api/pipelines/stable_diffusion/inpaint
title: Inpaint
- local: api/pipelines/stable_diffusion/depth2img
title: Depth-to-Image
- local: api/pipelines/stable_diffusion/image_variation
title: Image-Variation
- local: api/pipelines/stable_diffusion/upscale
title: Super-Resolution
- local: api/pipelines/stable_diffusion/latent_upscale
title: Stable-Diffusion-Latent-Upscaler
- local: api/pipelines/stable_diffusion/pix2pix
title: InstructPix2Pix
title: Stable Diffusion
- local: api/pipelines/stable_diffusion_2
title: Stable Diffusion 2
- local: api/pipelines/stochastic_karras_ve
title: Stochastic Karras VE
- local: api/pipelines/unclip
title: UnCLIP
- local: api/pipelines/latent_diffusion_uncond
title: Unconditional Latent Diffusion
- local: api/pipelines/versatile_diffusion
title: Versatile Diffusion
- local: api/pipelines/vq_diffusion
title: VQ Diffusion
title: Pipelines
- sections:
- local: api/schedulers/overview
title: Overview
- local: api/schedulers/ddim
title: DDIM
- local: api/schedulers/ddpm
title: DDPM
- local: api/schedulers/deis
title: DEIS
- local: api/schedulers/dpm_discrete
title: DPM Discrete Scheduler
- local: api/schedulers/dpm_discrete_ancestral
title: DPM Discrete Scheduler with ancestral sampling
- local: api/schedulers/euler_ancestral
title: Euler Ancestral Scheduler
- local: api/schedulers/euler
title: Euler scheduler
- local: api/schedulers/heun
title: Heun Scheduler
- local: api/schedulers/ipndm
title: IPNDM
- local: api/schedulers/lms_discrete
title: Linear Multistep
- local: api/schedulers/multistep_dpm_solver
title: Multistep DPM-Solver
- local: api/schedulers/pndm
title: PNDM
- local: api/schedulers/repaint
title: RePaint Scheduler
- local: api/schedulers/singlestep_dpm_solver
title: Singlestep DPM-Solver
- local: api/schedulers/stochastic_karras_ve
title: Stochastic Kerras VE
- local: api/schedulers/score_sde_ve
title: VE-SDE
- local: api/schedulers/score_sde_vp
title: VP-SDE
- local: api/schedulers/vq_diffusion
title: VQDiffusionScheduler
title: Schedulers
- sections:
- local: api/experimental/rl
title: RL Planning
title: Experimental Features
title: API

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Loaders
There are many ways to train adapter neural networks for diffusion models, such as
- [Textual Inversion](./training/text_inversion.mdx)
- [LoRA](https://github.com/cloneofsimo/lora)
- [Hypernetworks](https://arxiv.org/abs/1609.09106)
Such adapter neural networks often only consist of a fraction of the number of weights compared
to the pretrained model and as such are very portable. The Diffusers library offers an easy-to-use
API to load such adapter neural networks via the [`loaders.py` module](https://github.com/huggingface/diffusers/blob/main/src/diffusers/loaders.py).
**Note**: This module is still highly experimental and prone to future changes.
## LoaderMixins
### UNet2DConditionLoadersMixin
[[autodoc]] loaders.UNet2DConditionLoadersMixin

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Audio Diffusion
## Overview
[Audio Diffusion](https://github.com/teticio/audio-diffusion) by Robert Dargavel Smith.
Audio Diffusion leverages the recent advances in image generation using diffusion models by converting audio samples to
and from mel spectrogram images.
The original codebase of this implementation can be found [here](https://github.com/teticio/audio-diffusion), including
training scripts and example notebooks.
## Available Pipelines:
| Pipeline | Tasks | Colab
|---|---|:---:|
| [pipeline_audio_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/audio_diffusion/pipeline_audio_diffusion.py) | *Unconditional Audio Generation* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/teticio/audio-diffusion/blob/master/notebooks/audio_diffusion_pipeline.ipynb) |
## Examples:
### Audio Diffusion
```python
import torch
from IPython.display import Audio
from diffusers import DiffusionPipeline
device = "cuda" if torch.cuda.is_available() else "cpu"
pipe = DiffusionPipeline.from_pretrained("teticio/audio-diffusion-256").to(device)
output = pipe()
display(output.images[0])
display(Audio(output.audios[0], rate=mel.get_sample_rate()))
```
### Latent Audio Diffusion
```python
import torch
from IPython.display import Audio
from diffusers import DiffusionPipeline
device = "cuda" if torch.cuda.is_available() else "cpu"
pipe = DiffusionPipeline.from_pretrained("teticio/latent-audio-diffusion-256").to(device)
output = pipe()
display(output.images[0])
display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
```
### Audio Diffusion with DDIM (faster)
```python
import torch
from IPython.display import Audio
from diffusers import DiffusionPipeline
device = "cuda" if torch.cuda.is_available() else "cpu"
pipe = DiffusionPipeline.from_pretrained("teticio/audio-diffusion-ddim-256").to(device)
output = pipe()
display(output.images[0])
display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
```
### Variations, in-painting, out-painting etc.
```python
output = pipe(
raw_audio=output.audios[0, 0],
start_step=int(pipe.get_default_steps() / 2),
mask_start_secs=1,
mask_end_secs=1,
)
display(output.images[0])
display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
```
## AudioDiffusionPipeline
[[autodoc]] AudioDiffusionPipeline
- all
- __call__
## Mel
[[autodoc]] Mel

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Scalable Diffusion Models with Transformers (DiT)
## Overview
[Scalable Diffusion Models with Transformers](https://arxiv.org/abs/2212.09748) (DiT) by William Peebles and Saining Xie.
The abstract of the paper is the following:
*We explore a new class of diffusion models based on the transformer architecture. We train latent diffusion models of images, replacing the commonly-used U-Net backbone with a transformer that operates on latent patches. We analyze the scalability of our Diffusion Transformers (DiTs) through the lens of forward pass complexity as measured by Gflops. We find that DiTs with higher Gflops -- through increased transformer depth/width or increased number of input tokens -- consistently have lower FID. In addition to possessing good scalability properties, our largest DiT-XL/2 models outperform all prior diffusion models on the class-conditional ImageNet 512x512 and 256x256 benchmarks, achieving a state-of-the-art FID of 2.27 on the latter.*
The original codebase of this paper can be found here: [facebookresearch/dit](https://github.com/facebookresearch/dit).
## Available Pipelines:
| Pipeline | Tasks | Colab
|---|---|:---:|
| [pipeline_dit.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/dit/pipeline_dit.py) | *Conditional Image Generation* | - |
## Usage example
```python
from diffusers import DiTPipeline, DPMSolverMultistepScheduler
import torch
pipe = DiTPipeline.from_pretrained("facebook/DiT-XL-2-256", torch_dtype=torch.float16)
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
pipe = pipe.to("cuda")
# pick words from Imagenet class labels
pipe.labels # to print all available words
# pick words that exist in ImageNet
words = ["white shark", "umbrella"]
class_ids = pipe.get_label_ids(words)
generator = torch.manual_seed(33)
output = pipe(class_labels=class_ids, num_inference_steps=25, generator=generator)
image = output.images[0] # label 'white shark'
```
## DiTPipeline
[[autodoc]] DiTPipeline
- all
- __call__

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# PaintByExample
## Overview
[Paint by Example: Exemplar-based Image Editing with Diffusion Models](https://arxiv.org/abs/2211.13227) by Binxin Yang, Shuyang Gu, Bo Zhang, Ting Zhang, Xuejin Chen, Xiaoyan Sun, Dong Chen, Fang Wen
The abstract of the paper is the following:
*Language-guided image editing has achieved great success recently. In this paper, for the first time, we investigate exemplar-guided image editing for more precise control. We achieve this goal by leveraging self-supervised training to disentangle and re-organize the source image and the exemplar. However, the naive approach will cause obvious fusing artifacts. We carefully analyze it and propose an information bottleneck and strong augmentations to avoid the trivial solution of directly copying and pasting the exemplar image. Meanwhile, to ensure the controllability of the editing process, we design an arbitrary shape mask for the exemplar image and leverage the classifier-free guidance to increase the similarity to the exemplar image. The whole framework involves a single forward of the diffusion model without any iterative optimization. We demonstrate that our method achieves an impressive performance and enables controllable editing on in-the-wild images with high fidelity.*
The original codebase can be found [here](https://github.com/Fantasy-Studio/Paint-by-Example).
## Available Pipelines:
| Pipeline | Tasks | Colab
|---|---|:---:|
| [pipeline_paint_by_example.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/paint_by_example/pipeline_paint_by_example.py) | *Image-Guided Image Painting* | - |
## Tips
- PaintByExample is supported by the official [Fantasy-Studio/Paint-by-Example](https://huggingface.co/Fantasy-Studio/Paint-by-Example) checkpoint. The checkpoint has been warm-started from the [CompVis/stable-diffusion-v1-4](https://huggingface.co/CompVis/stable-diffusion-v1-4) and with the objective to inpaint partly masked images conditioned on example / reference images
- To quickly demo *PaintByExample*, please have a look at [this demo](https://huggingface.co/spaces/Fantasy-Studio/Paint-by-Example)
- You can run the following code snippet as an example:
```python
# !pip install diffusers transformers
import PIL
import requests
import torch
from io import BytesIO
from diffusers import DiffusionPipeline
def download_image(url):
response = requests.get(url)
return PIL.Image.open(BytesIO(response.content)).convert("RGB")
img_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/image/example_1.png"
mask_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/mask/example_1.png"
example_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/reference/example_1.jpg"
init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
example_image = download_image(example_url).resize((512, 512))
pipe = DiffusionPipeline.from_pretrained(
"Fantasy-Studio/Paint-by-Example",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
image = pipe(image=init_image, mask_image=mask_image, example_image=example_image).images[0]
image
```
## PaintByExamplePipeline
[[autodoc]] PaintByExamplePipeline
- all
- __call__

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Depth-to-Image Generation
## StableDiffusionDepth2ImgPipeline
The depth-guided stable diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), and [LAION](https://laion.ai/), as part of Stable Diffusion 2.0. It uses [MiDas](https://github.com/isl-org/MiDaS) to infer depth based on an image.
[`StableDiffusionDepth2ImgPipeline`] lets you pass a text prompt and an initial image to condition the generation of new images as well as a `depth_map` to preserve the images structure.
The original codebase can be found here:
- *Stable Diffusion v2*: [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion#depth-conditional-stable-diffusion)
Available Checkpoints are:
- *stable-diffusion-2-depth*: [stabilityai/stable-diffusion-2-depth](https://huggingface.co/stabilityai/stable-diffusion-2-depth)
[[autodoc]] StableDiffusionDepth2ImgPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Image Variation
## StableDiffusionImageVariationPipeline
[`StableDiffusionImageVariationPipeline`] lets you generate variations from an input image using Stable Diffusion. It uses a fine-tuned version of Stable Diffusion model, trained by [Justin Pinkney](https://www.justinpinkney.com/) (@Buntworthy) at [Lambda](https://lambdalabs.com/)
The original codebase can be found here:
[Stable Diffusion Image Variations](https://github.com/LambdaLabsML/lambda-diffusers#stable-diffusion-image-variations)
Available Checkpoints are:
- *sd-image-variations-diffusers*: [lambdalabs/sd-image-variations-diffusers](https://huggingface.co/lambdalabs/sd-image-variations-diffusers)
[[autodoc]] StableDiffusionImageVariationPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Image-to-Image Generation
## StableDiffusionImg2ImgPipeline
The Stable Diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionImg2ImgPipeline`] lets you pass a text prompt and an initial image to condition the generation of new images using Stable Diffusion.
The original codebase can be found here: [CampVis/stable-diffusion](https://github.com/CompVis/stable-diffusion/blob/main/scripts/img2img.py)
[`StableDiffusionImg2ImgPipeline`] is compatible with all Stable Diffusion checkpoints for [Text-to-Image](./text2img)
[[autodoc]] StableDiffusionImg2ImgPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Text-Guided Image Inpainting
## StableDiffusionInpaintPipeline
The Stable Diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionInpaintPipeline`] lets you edit specific parts of an image by providing a mask and a text prompt using Stable Diffusion.
The original codebase can be found here:
- *Stable Diffusion V1*: [CampVis/stable-diffusion](https://github.com/runwayml/stable-diffusion#inpainting-with-stable-diffusion)
- *Stable Diffusion V2*: [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion#image-inpainting-with-stable-diffusion)
Available checkpoints are:
- *stable-diffusion-inpainting (512x512 resolution)*: [runwayml/stable-diffusion-inpainting](https://huggingface.co/runwayml/stable-diffusion-inpainting)
- *stable-diffusion-2-inpainting (512x512 resolution)*: [stabilityai/stable-diffusion-2-inpainting](https://huggingface.co/stabilityai/stable-diffusion-2-inpainting)
[[autodoc]] StableDiffusionInpaintPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention

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<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Stable Diffusion Latent Upscaler
## StableDiffusionLatentUpscalePipeline
The Stable Diffusion Latent Upscaler model was created by [Katherine Crowson](https://github.com/crowsonkb/k-diffusion) in collaboration with [Stability AI](https://stability.ai/). It can be used on top of any [`StableDiffusionUpscalePipeline`] checkpoint to enhance its output image resolution by a factor of 2.
A notebook that demonstrates the original implementation can be found here:
- [Stable Diffusion Upscaler Demo](https://colab.research.google.com/drive/1o1qYJcFeywzCIdkfKJy7cTpgZTCM2EI4)
Available Checkpoints are:
- *stabilityai/latent-upscaler*: [stabilityai/sd-x2-latent-upscaler](https://huggingface.co/stabilityai/sd-x2-latent-upscaler)
[[autodoc]] StableDiffusionLatentUpscalePipeline
- all
- __call__
- enable_sequential_cpu_offload
- enable_attention_slicing
- disable_attention_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention

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<!--Copyright 2023 The HuggingFace Team. All rights reserved.
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# InstructPix2Pix: Learning to Follow Image Editing Instructions
## Overview
[InstructPix2Pix: Learning to Follow Image Editing Instructions](https://arxiv.org/abs/2211.09800) by Tim Brooks, Aleksander Holynski and Alexei A. Efros.
The abstract of the paper is the following:
*We propose a method for editing images from human instructions: given an input image and a written instruction that tells the model what to do, our model follows these instructions to edit the image. To obtain training data for this problem, we combine the knowledge of two large pretrained models -- a language model (GPT-3) and a text-to-image model (Stable Diffusion) -- to generate a large dataset of image editing examples. Our conditional diffusion model, InstructPix2Pix, is trained on our generated data, and generalizes to real images and user-written instructions at inference time. Since it performs edits in the forward pass and does not require per example fine-tuning or inversion, our model edits images quickly, in a matter of seconds. We show compelling editing results for a diverse collection of input images and written instructions.*
Resources:
* [Project Page](https://www.timothybrooks.com/instruct-pix2pix).
* [Paper](https://arxiv.org/abs/2211.09800).
* [Original Code](https://github.com/timothybrooks/instruct-pix2pix).
* [Demo](https://huggingface.co/spaces/timbrooks/instruct-pix2pix).
## Available Pipelines:
| Pipeline | Tasks | Demo
|---|---|:---:|
| [StableDiffusionInstructPix2PixPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_instruct_pix2pix.py) | *Text-Based Image Editing* | [🤗 Space](https://huggingface.co/spaces/timbrooks/instruct-pix2pix) |
<!-- TODO: add Colab -->
## Usage example
```python
import PIL
import requests
import torch
from diffusers import StableDiffusionInstructPix2PixPipeline
model_id = "timbrooks/instruct-pix2pix"
pipe = StableDiffusionInstructPix2PixPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
url = "https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/mountain.png"
def download_image(url):
image = PIL.Image.open(requests.get(url, stream=True).raw)
image = PIL.ImageOps.exif_transpose(image)
image = image.convert("RGB")
return image
image = download_image(url)
prompt = "make the mountains snowy"
edit = pipe(prompt, image=image, num_inference_steps=20, image_guidance_scale=1.5, guidance_scale=7).images[0]
images[0].save("snowy_mountains.png")
```
## StableDiffusionInstructPix2PixPipeline
[[autodoc]] StableDiffusionInstructPix2PixPipeline
- __call__
- all

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Text-to-Image Generation
## StableDiffusionPipeline
The Stable Diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionPipeline`] is capable of generating photo-realistic images given any text input using Stable Diffusion.
The original codebase can be found here:
- *Stable Diffusion V1*: [CampVis/stable-diffusion](https://github.com/CompVis/stable-diffusion)
- *Stable Diffusion v2*: [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion)
Available Checkpoints are:
- *stable-diffusion-v1-4 (512x512 resolution)* [CompVis/stable-diffusion-v1-4](https://huggingface.co/CompVis/stable-diffusion-v1-4)
- *stable-diffusion-v1-5 (512x512 resolution)* [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5)
- *stable-diffusion-2-base (512x512 resolution)*: [stabilityai/stable-diffusion-2-base](https://huggingface.co/stabilityai/stable-diffusion-2-base)
- *stable-diffusion-2 (768x768 resolution)*: [stabilityai/stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2)
- *stable-diffusion-2-1-base (512x512 resolution)* [stabilityai/stable-diffusion-2-1-base](https://huggingface.co/stabilityai/stable-diffusion-2-1-base)
- *stable-diffusion-2-1 (768x768 resolution)*: [stabilityai/stable-diffusion-2-1](https://huggingface.co/stabilityai/stable-diffusion-2-1)
[[autodoc]] StableDiffusionPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_vae_slicing
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Super-Resolution
## StableDiffusionUpscalePipeline
The upscaler diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), and [LAION](https://laion.ai/), as part of Stable Diffusion 2.0. [`StableDiffusionUpscalePipeline`] can be used to enhance the resolution of input images by a factor of 4.
The original codebase can be found here:
- *Stable Diffusion v2*: [Stability-AI/stablediffusion](https://github.com/Stability-AI/stablediffusion#image-upscaling-with-stable-diffusion)
Available Checkpoints are:
- *stabilityai/stable-diffusion-x4-upscaler (x4 resolution resolution)*: [stable-diffusion-x4-upscaler](https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler)
[[autodoc]] StableDiffusionUpscalePipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Stable diffusion 2
Stable Diffusion 2 is a text-to-image _latent diffusion_ model built upon the work of [Stable Diffusion 1](https://stability.ai/blog/stable-diffusion-public-release).
The project to train Stable Diffusion 2 was led by Robin Rombach and Katherine Crowson from [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/).
*The Stable Diffusion 2.0 release includes robust text-to-image models trained using a brand new text encoder (OpenCLIP), developed by LAION with support from Stability AI, which greatly improves the quality of the generated images compared to earlier V1 releases. The text-to-image models in this release can generate images with default resolutions of both 512x512 pixels and 768x768 pixels.
These models are trained on an aesthetic subset of the [LAION-5B dataset](https://laion.ai/blog/laion-5b/) created by the DeepFloyd team at Stability AI, which is then further filtered to remove adult content using [LAIONs NSFW filter](https://openreview.net/forum?id=M3Y74vmsMcY).*
For more details about how Stable Diffusion 2 works and how it differs from Stable Diffusion 1, please refer to the official [launch announcement post](https://stability.ai/blog/stable-diffusion-v2-release).
## Tips
### Available checkpoints:
Note that the architecture is more or less identical to [Stable Diffusion 1](./stable_diffusion/overview) so please refer to [this page](./stable_diffusion/overview) for API documentation.
- *Text-to-Image (512x512 resolution)*: [stabilityai/stable-diffusion-2-base](https://huggingface.co/stabilityai/stable-diffusion-2-base) with [`StableDiffusionPipeline`]
- *Text-to-Image (768x768 resolution)*: [stabilityai/stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2) with [`StableDiffusionPipeline`]
- *Image Inpainting (512x512 resolution)*: [stabilityai/stable-diffusion-2-inpainting](https://huggingface.co/stabilityai/stable-diffusion-2-inpainting) with [`StableDiffusionInpaintPipeline`]
- *Super-Resolution (x4 resolution resolution)*: [stable-diffusion-x4-upscaler](https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler) [`StableDiffusionUpscalePipeline`]
- *Depth-to-Image (512x512 resolution)*: [stabilityai/stable-diffusion-2-depth](https://huggingface.co/stabilityai/stable-diffusion-2-depth) with [`StableDiffusionDepth2ImagePipeline`]
We recommend using the [`DPMSolverMultistepScheduler`] as it's currently the fastest scheduler there is.
### Text-to-Image
- *Text-to-Image (512x512 resolution)*: [stabilityai/stable-diffusion-2-base](https://huggingface.co/stabilityai/stable-diffusion-2-base) with [`StableDiffusionPipeline`]
```python
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
import torch
repo_id = "stabilityai/stable-diffusion-2-base"
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
pipe = pipe.to("cuda")
prompt = "High quality photo of an astronaut riding a horse in space"
image = pipe(prompt, num_inference_steps=25).images[0]
image.save("astronaut.png")
```
- *Text-to-Image (768x768 resolution)*: [stabilityai/stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2) with [`StableDiffusionPipeline`]
```python
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
import torch
repo_id = "stabilityai/stable-diffusion-2"
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
pipe = pipe.to("cuda")
prompt = "High quality photo of an astronaut riding a horse in space"
image = pipe(prompt, guidance_scale=9, num_inference_steps=25).images[0]
image.save("astronaut.png")
```
### Image Inpainting
- *Image Inpainting (512x512 resolution)*: [stabilityai/stable-diffusion-2-inpainting](https://huggingface.co/stabilityai/stable-diffusion-2-inpainting) with [`StableDiffusionInpaintPipeline`]
```python
import PIL
import requests
import torch
from io import BytesIO
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
def download_image(url):
response = requests.get(url)
return PIL.Image.open(BytesIO(response.content)).convert("RGB")
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
repo_id = "stabilityai/stable-diffusion-2-inpainting"
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
pipe = pipe.to("cuda")
prompt = "Face of a yellow cat, high resolution, sitting on a park bench"
image = pipe(prompt=prompt, image=init_image, mask_image=mask_image, num_inference_steps=25).images[0]
image.save("yellow_cat.png")
```
### Super-Resolution
- *Image Upscaling (x4 resolution resolution)*: [stable-diffusion-x4-upscaler](https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler) with [`StableDiffusionUpscalePipeline`]
```python
import requests
from PIL import Image
from io import BytesIO
from diffusers import StableDiffusionUpscalePipeline
import torch
# load model and scheduler
model_id = "stabilityai/stable-diffusion-x4-upscaler"
pipeline = StableDiffusionUpscalePipeline.from_pretrained(model_id, torch_dtype=torch.float16)
pipeline = pipeline.to("cuda")
# let's download an image
url = "https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd2-upscale/low_res_cat.png"
response = requests.get(url)
low_res_img = Image.open(BytesIO(response.content)).convert("RGB")
low_res_img = low_res_img.resize((128, 128))
prompt = "a white cat"
upscaled_image = pipeline(prompt=prompt, image=low_res_img).images[0]
upscaled_image.save("upsampled_cat.png")
```
### Depth-to-Image
- *Depth-Guided Text-to-Image*: [stabilityai/stable-diffusion-2-depth](https://huggingface.co/stabilityai/stable-diffusion-2-depth) [`StableDiffusionDepth2ImagePipeline`]
```python
import torch
import requests
from PIL import Image
from diffusers import StableDiffusionDepth2ImgPipeline
pipe = StableDiffusionDepth2ImgPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-depth",
torch_dtype=torch.float16,
).to("cuda")
url = "http://images.cocodataset.org/val2017/000000039769.jpg"
init_image = Image.open(requests.get(url, stream=True).raw)
prompt = "two tigers"
n_propmt = "bad, deformed, ugly, bad anotomy"
image = pipe(prompt=prompt, image=init_image, negative_prompt=n_propmt, strength=0.7).images[0]
```
### How to load and use different schedulers.
The stable diffusion pipeline uses [`DDIMScheduler`] scheduler by default. But `diffusers` provides many other schedulers that can be used with the stable diffusion pipeline such as [`PNDMScheduler`], [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`] etc.
To use a different scheduler, you can either change it via the [`ConfigMixin.from_config`] method or pass the `scheduler` argument to the `from_pretrained` method of the pipeline. For example, to use the [`EulerDiscreteScheduler`], you can do the following:
```python
>>> from diffusers import StableDiffusionPipeline, EulerDiscreteScheduler
>>> pipeline = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2")
>>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
>>> # or
>>> euler_scheduler = EulerDiscreteScheduler.from_pretrained("stabilityai/stable-diffusion-2", subfolder="scheduler")
>>> pipeline = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2", scheduler=euler_scheduler)
```

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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# unCLIP
## Overview
[Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) by Aditya Ramesh, Prafulla Dhariwal, Alex Nichol, Casey Chu, Mark Chen
The abstract of the paper is the following:
Contrastive models like CLIP have been shown to learn robust representations of images that capture both semantics and style. To leverage these representations for image generation, we propose a two-stage model: a prior that generates a CLIP image embedding given a text caption, and a decoder that generates an image conditioned on the image embedding. We show that explicitly generating image representations improves image diversity with minimal loss in photorealism and caption similarity. Our decoders conditioned on image representations can also produce variations of an image that preserve both its semantics and style, while varying the non-essential details absent from the image representation. Moreover, the joint embedding space of CLIP enables language-guided image manipulations in a zero-shot fashion. We use diffusion models for the decoder and experiment with both autoregressive and diffusion models for the prior, finding that the latter are computationally more efficient and produce higher-quality samples.
The unCLIP model in diffusers comes from kakaobrain's karlo and the original codebase can be found [here](https://github.com/kakaobrain/karlo). Additionally, lucidrains has a DALL-E 2 recreation [here](https://github.com/lucidrains/DALLE2-pytorch).
## Available Pipelines:
| Pipeline | Tasks | Colab
|---|---|:---:|
| [pipeline_unclip.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/unclip/pipeline_unclip.py) | *Text-to-Image Generation* | - |
| [pipeline_unclip_image_variation.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/unclip/pipeline_unclip_image_variation.py) | *Image-Guided Image Generation* | - |
## UnCLIPPipeline
[[autodoc]] UnCLIPPipeline
- all
- __call__
[[autodoc]] UnCLIPImageVariationPipeline
- all
- __call__

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Denoising diffusion implicit models (DDIM)
## Overview
[Denoising Diffusion Implicit Models](https://arxiv.org/abs/2010.02502) (DDIM) by Jiaming Song, Chenlin Meng and Stefano Ermon.
The abstract of the paper is the following:
Denoising diffusion probabilistic models (DDPMs) have achieved high quality image generation without adversarial training, yet they require simulating a Markov chain for many steps to produce a sample. To accelerate sampling, we present denoising diffusion implicit models (DDIMs), a more efficient class of iterative implicit probabilistic models with the same training procedure as DDPMs. In DDPMs, the generative process is defined as the reverse of a Markovian diffusion process. We construct a class of non-Markovian diffusion processes that lead to the same training objective, but whose reverse process can be much faster to sample from. We empirically demonstrate that DDIMs can produce high quality samples 10× to 50× faster in terms of wall-clock time compared to DDPMs, allow us to trade off computation for sample quality, and can perform semantically meaningful image interpolation directly in the latent space.
The original codebase of this paper can be found here: [ermongroup/ddim](https://github.com/ermongroup/ddim).
For questions, feel free to contact the author on [tsong.me](https://tsong.me/).
## DDIMScheduler
[[autodoc]] DDIMScheduler

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Denoising diffusion probabilistic models (DDPM)
## Overview
[Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2006.11239)
(DDPM) by Jonathan Ho, Ajay Jain and Pieter Abbeel proposes the diffusion based model of the same name, but in the context of the 🤗 Diffusers library, DDPM refers to the discrete denoising scheduler from the paper as well as the pipeline.
The abstract of the paper is the following:
We present high quality image synthesis results using diffusion probabilistic models, a class of latent variable models inspired by considerations from nonequilibrium thermodynamics. Our best results are obtained by training on a weighted variational bound designed according to a novel connection between diffusion probabilistic models and denoising score matching with Langevin dynamics, and our models naturally admit a progressive lossy decompression scheme that can be interpreted as a generalization of autoregressive decoding. On the unconditional CIFAR10 dataset, we obtain an Inception score of 9.46 and a state-of-the-art FID score of 3.17. On 256x256 LSUN, we obtain sample quality similar to ProgressiveGAN.
The original paper can be found [here](https://arxiv.org/abs/2010.02502).
## DDPMScheduler
[[autodoc]] DDPMScheduler

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# DEIS
Fast Sampling of Diffusion Models with Exponential Integrator.
## Overview
Original paper can be found [here](https://arxiv.org/abs/2204.13902). The original implementation can be found [here](https://github.com/qsh-zh/deis).
## DEISMultistepScheduler
[[autodoc]] DEISMultistepScheduler

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# DPM Discrete Scheduler inspired by Karras et. al paper
## Overview
Inspired by [Karras et. al](https://arxiv.org/abs/2206.00364). Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
## KDPM2DiscreteScheduler
[[autodoc]] KDPM2DiscreteScheduler

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# DPM Discrete Scheduler with ancestral sampling inspired by Karras et. al paper
## Overview
Inspired by [Karras et. al](https://arxiv.org/abs/2206.00364). Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
## KDPM2AncestralDiscreteScheduler
[[autodoc]] KDPM2AncestralDiscreteScheduler

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Euler scheduler
## Overview
Euler scheduler (Algorithm 2) from the paper [Elucidating the Design Space of Diffusion-Based Generative Models](https://arxiv.org/abs/2206.00364) by Karras et al. (2022). Based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L51) implementation by Katherine Crowson.
Fast scheduler which often times generates good outputs with 20-30 steps.
## EulerDiscreteScheduler
[[autodoc]] EulerDiscreteScheduler

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Euler Ancestral scheduler
## Overview
Ancestral sampling with Euler method steps. Based on the original (k-diffusion)[https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72] implementation by Katherine Crowson.
Fast scheduler which often times generates good outputs with 20-30 steps.
## EulerAncestralDiscreteScheduler
[[autodoc]] EulerAncestralDiscreteScheduler

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Heun scheduler inspired by Karras et. al paper
## Overview
Algorithm 1 of [Karras et. al](https://arxiv.org/abs/2206.00364).
Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
## HeunDiscreteScheduler
[[autodoc]] HeunDiscreteScheduler

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# improved pseudo numerical methods for diffusion models (iPNDM)
## Overview
Original implementation can be found [here](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296).
## IPNDMScheduler
[[autodoc]] IPNDMScheduler

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specific language governing permissions and limitations under the License.
-->
# Linear multistep scheduler for discrete beta schedules
## Overview
Original implementation can be found [here](https://arxiv.org/abs/2206.00364).
## LMSDiscreteScheduler
[[autodoc]] LMSDiscreteScheduler

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# Multistep DPM-Solver
## Overview
Original paper can be found [here](https://arxiv.org/abs/2206.00927) and the [improved version](https://arxiv.org/abs/2211.01095). The original implementation can be found [here](https://github.com/LuChengTHU/dpm-solver).
## DPMSolverMultistepScheduler
[[autodoc]] DPMSolverMultistepScheduler

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# Schedulers
Diffusers contains multiple pre-built schedule functions for the diffusion process.
## What is a scheduler?
The schedule functions, denoted *Schedulers* in the library take in the output of a trained model, a sample which the diffusion process is iterating on, and a timestep to return a denoised sample. That's why schedulers may also be called *Samplers* in other diffusion models implementations.
- Schedulers define the methodology for iteratively adding noise to an image or for updating a sample based on model outputs.
- adding noise in different manners represent the algorithmic processes to train a diffusion model by adding noise to images.
- for inference, the scheduler defines how to update a sample based on an output from a pretrained model.
- Schedulers are often defined by a *noise schedule* and an *update rule* to solve the differential equation solution.
### Discrete versus continuous schedulers
All schedulers take in a timestep to predict the updated version of the sample being diffused.
The timesteps dictate where in the diffusion process the step is, where data is generated by iterating forward in time and inference is executed by propagating backwards through timesteps.
Different algorithms use timesteps that can be discrete (accepting `int` inputs), such as the [`DDPMScheduler`] or [`PNDMScheduler`], or continuous (accepting `float` inputs), such as the score-based schedulers [`ScoreSdeVeScheduler`] or [`ScoreSdeVpScheduler`].
## Designing Re-usable schedulers
The core design principle between the schedule functions is to be model, system, and framework independent.
This allows for rapid experimentation and cleaner abstractions in the code, where the model prediction is separated from the sample update.
To this end, the design of schedulers is such that:
- Schedulers can be used interchangeably between diffusion models in inference to find the preferred trade-off between speed and generation quality.
- Schedulers are currently by default in PyTorch, but are designed to be framework independent (partial Jax support currently exists).
- Many diffusion pipelines, such as [`StableDiffusionPipeline`] and [`DiTPipeline`] can use any of [`KarrasDiffusionSchedulers`]
## Schedulers Summary
The following table summarizes all officially supported schedulers, their corresponding paper
| Scheduler | Paper |
|---|---|
| [ddim](./ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) |
| [ddpm](./ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) |
| [singlestep_dpm_solver](./singlestep_dpm_solver) | [**Singlestep DPM-Solver**](https://arxiv.org/abs/2206.00927) |
| [multistep_dpm_solver](./multistep_dpm_solver) | [**Multistep DPM-Solver**](https://arxiv.org/abs/2206.00927) |
| [heun](./heun) | [**Heun scheduler inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
| [dpm_discrete](./dpm_discrete) | [**DPM Discrete Scheduler inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
| [dpm_discrete_ancestral](./dpm_discrete_ancestral) | [**DPM Discrete Scheduler with ancestral sampling inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
| [stochastic_karras_ve](./stochastic_karras_ve) | [**Variance exploding, stochastic sampling from Karras et. al**](https://arxiv.org/abs/2206.00364) |
| [lms_discrete](./lms_discrete) | [**Linear multistep scheduler for discrete beta schedules**](https://arxiv.org/abs/2206.00364) |
| [pndm](./pndm) | [**Pseudo numerical methods for diffusion models (PNDM)**](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181) |
| [score_sde_ve](./score_sde_ve) | [**variance exploding stochastic differential equation (VE-SDE) scheduler**](https://arxiv.org/abs/2011.13456) |
| [ipndm](./ipndm) | [**improved pseudo numerical methods for diffusion models (iPNDM)**](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296) |
| [score_sde_vp](./score_sde_vp) | [**Variance preserving stochastic differential equation (VP-SDE) scheduler**](https://arxiv.org/abs/2011.13456) |
| [euler](./euler) | [**Euler scheduler**](https://arxiv.org/abs/2206.00364) |
| [euler_ancestral](./euler_ancestral) | [**Euler Ancestral scheduler**](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72) |
| [vq_diffusion](./vq_diffusion) | [**VQDiffusionScheduler**](https://arxiv.org/abs/2111.14822) |
| [repaint](./repaint) | [**RePaint scheduler**](https://arxiv.org/abs/2201.09865) |
## API
The core API for any new scheduler must follow a limited structure.
- Schedulers should provide one or more `def step(...)` functions that should be called to update the generated sample iteratively.
- Schedulers should provide a `set_timesteps(...)` method that configures the parameters of a schedule function for a specific inference task.
- Schedulers should be framework-specific.
The base class [`SchedulerMixin`] implements low level utilities used by multiple schedulers.
### SchedulerMixin
[[autodoc]] SchedulerMixin
### SchedulerOutput
The class [`SchedulerOutput`] contains the outputs from any schedulers `step(...)` call.
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
### KarrasDiffusionSchedulers
[[autodoc]] schedulers.scheduling_utils.KarrasDiffusionSchedulers

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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
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# Pseudo numerical methods for diffusion models (PNDM)
## Overview
Original implementation can be found [here](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181).
## PNDMScheduler
[[autodoc]] PNDMScheduler

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# RePaint scheduler
## Overview
DDPM-based inpainting scheduler for unsupervised inpainting with extreme masks.
Intended for use with [`RePaintPipeline`].
Based on the paper [RePaint: Inpainting using Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2201.09865)
and the original implementation by Andreas Lugmayr et al.: https://github.com/andreas128/RePaint
## RePaintScheduler
[[autodoc]] RePaintScheduler

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an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# variance exploding stochastic differential equation (VE-SDE) scheduler
## Overview
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
## ScoreSdeVeScheduler
[[autodoc]] ScoreSdeVeScheduler

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an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
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-->
# Variance preserving stochastic differential equation (VP-SDE) scheduler
## Overview
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
<Tip warning={true}>
Score SDE-VP is under construction.
</Tip>
## ScoreSdeVpScheduler
[[autodoc]] schedulers.scheduling_sde_vp.ScoreSdeVpScheduler

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# Singlestep DPM-Solver
## Overview
Original paper can be found [here](https://arxiv.org/abs/2206.00927) and the [improved version](https://arxiv.org/abs/2211.01095). The original implementation can be found [here](https://github.com/LuChengTHU/dpm-solver).
## DPMSolverSinglestepScheduler
[[autodoc]] DPMSolverSinglestepScheduler

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an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Variance exploding, stochastic sampling from Karras et. al
## Overview
Original paper can be found [here](https://arxiv.org/abs/2206.00364).
## KarrasVeScheduler
[[autodoc]] KarrasVeScheduler

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specific language governing permissions and limitations under the License.
-->
# VQDiffusionScheduler
## Overview
Original paper can be found [here](https://arxiv.org/abs/2111.14822)
## VQDiffusionScheduler
[[autodoc]] VQDiffusionScheduler

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# Philosophy
🧨 Diffusers provides **state-of-the-art** pretrained diffusion models across multiple modalities.
Its purpose is to serve as a **modular toolbox** for both inference and training.
We aim at building a library that stands the test of time and therefore take API design very seriously.
In a nutshell, Diffusers is built to be a natural extension of PyTorch. Therefore, most of our design choices are based on [PyTorch's Design Principles](https://pytorch.org/docs/stable/community/design.html#pytorch-design-philosophy). Let's go over the most important ones:
## Usability over Performance
- While Diffusers has many built-in performance-enhancing features (see [Memory and Speed](https://huggingface.co/docs/diffusers/optimization/fp16)), models are always loaded with the highest precision and lowest optimization. Therefore, by default diffusion pipelines are always instantiated on CPU with float32 precision if not otherwise defined by the user. This ensures usability across different platforms and accelerators and means that no complex installations are required to run the library.
- Diffusers aim at being a **light-weight** package and therefore has very few required dependencies, but many soft dependencies that can improve performance (such as `accelerate`, `safetensors`, `onnx`, etc...). We strive to keep the library as lightweight as possible so that it can be added without much concern as a dependency on other packages.
- Diffusers prefers simple, self-explainable code over condensed, magic code. This means that short-hand code syntaxes such as lambda functions, and advanced PyTorch operators are often not desired.
## Simple over easy
As PyTorch states, **explicit is better than implicit** and **simple is better than complex**. This design philosophy is reflected in multiple parts of the library:
- We follow PyTorch's API with methods like [`DiffusionPipeline.to`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.to) to let the user handle device management.
- Raising concise error messages is preferred to silently correct erroneous input. Diffusers aims at teaching the user, rather than making the library as easy to use as possible.
- Complex model vs. scheduler logic is exposed instead of magically handled inside. Schedulers/Samplers are separated from diffusion models with minimal dependencies on each other. This forces the user to write the unrolled denoising loop. However, the separation allows for easier debugging and gives the user more control over adapting the denoising process or switching out diffusion models or schedulers.
- Separately trained components of the diffusion pipeline, *e.g.* the text encoder, the unet, and the variational autoencoder, each have their own model class. This forces the user to handle the interaction between the different model components, and the serialization format separates the model components into different files. However, this allows for easier debugging and customization. Dreambooth or textual inversion training
is very simple thanks to diffusers' ability to separate single components of the diffusion pipeline.
## Tweakable, contributor-friendly over abstraction
For large parts of the library, Diffusers adopts an important design principle of the [Transformers library](https://github.com/huggingface/transformers), which is to prefer copy-pasted code over hasty abstractions. This design principle is very opinionated and stands in stark contrast to popular design principles such as [Don't repeat yourself (DRY)](https://en.wikipedia.org/wiki/Don%27t_repeat_yourself).
In short, just like Transformers does for modeling files, diffusers prefers to keep an extremely low level of abstraction and very self-contained code for pipelines and schedulers.
Functions, long code blocks, and even classes can be copied across multiple files which at first can look like a bad, sloppy design choice that makes the library unmaintainable.
**However**, this design has proven to be extremely successful for Transformers and makes a lot of sense for community-driven, open-source machine learning libraries because:
- Machine Learning is an extremely fast-moving field in which paradigms, model architectures, and algorithms are changing rapidly, which therefore makes it very difficult to define long-lasting code abstractions.
- Machine Learning practitioners like to be able to quickly tweak existing code for ideation and research and therefore prefer self-contained code over one that contains many abstractions.
- Open-source libraries rely on community contributions and therefore must build a library that is easy to contribute to. The more abstract the code, the more dependencies, the harder to read, and the harder to contribute to. Contributors simply stop contributing to very abstract libraries out of fear of breaking vital functionality. If contributing to a library cannot break other fundamental code, not only is it more inviting for potential new contributors, but it is also easier to review and contribute to multiple parts in parallel.
At Hugging Face, we call this design the **single-file policy** which means that almost all of the code of a certain class should be written in a single, self-contained file. To read more about the philosophy, you can have a look
at [this blog post](https://huggingface.co/blog/transformers-design-philosophy).
In diffusers, we follow this philosophy for both pipelines and schedulers, but only partly for diffusion models. The reason we don't follow this design fully for diffusion models is because almost all diffusion pipelines, such
as [DDPM](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/ddpm), [Stable Diffusion](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/stable_diffusion/overview#stable-diffusion-pipelines), [UnCLIP (Dalle-2)](https://huggingface.co/docs/diffusers/v0.12.0/en/api/pipelines/unclip#overview) and [Imagen](https://imagen.research.google/) all rely on the same diffusion model, the [UNet](https://huggingface.co/docs/diffusers/api/models#diffusers.UNet2DConditionModel).
Great, now you should have generally understood why 🧨 Diffusers is designed the way it is 🤗.
We try to apply these design principles consistently across the library. Nevertheless, there are some minor exceptions to the philosophy or some unlucky design choices. If you have feedback regarding the design, we would ❤️ to hear it [directly on GitHub](https://github.com/huggingface/diffusers/issues/new?assignees=&labels=&template=feedback.md&title=).
## Design Philosophy in Details
Now, let's look a bit into the nitty-gritty details of the design philosophy. Diffusers essentially consist of three major classes, [pipelines](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines), [models](https://github.com/huggingface/diffusers/tree/main/src/diffusers/models), and [schedulers](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers).
Let's walk through more in-detail design decisions for each class.
### Pipelines
Pipelines are designed to be easy to use (therefore do not follow [*Simple over easy*](#simple-over-easy) 100%)), are not feature complete, and should loosely be seen as examples of how to use [models](#models) and [schedulers](#schedulers) for inference.
The following design principles are followed:
- Pipelines follow the single-file policy. All pipelines can be found in individual directories under src/diffusers/pipelines. One pipeline folder corresponds to one diffusion paper/project/release. Multiple pipeline files can be gathered in one pipeline folder, as its done for [`src/diffusers/pipelines/stable-diffusion`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines/stable_diffusion). If pipelines share similar functionality, one can make use of the [#Copied from mechanism](https://github.com/huggingface/diffusers/blob/125d783076e5bd9785beb05367a2d2566843a271/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_img2img.py#L251).
- Pipelines all inherit from [`DiffusionPipeline`]
- Every pipeline consists of different model and scheduler components, that are documented in the [`model_index.json` file](https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/model_index.json), are accessible under the same name as attributes of the pipeline and can be shared between pipelines with [`DiffusionPipeline.components`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.components) function.
- Every pipeline should be loadable via the [`DiffusionPipeline.from_pretrained`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.from_pretrained) function.
- Pipelines should be used **only** for inference.
- Pipelines should be very readable, self-explanatory, and easy to tweak.
- Pipelines should be designed to build on top of each other and be easy to integrate into higher-level APIs.
- Pipelines are **not** intended to be feature-complete user interfaces. For future complete user interfaces one should rather have a look at [InvokeAI](https://github.com/invoke-ai/InvokeAI), [Diffuzers](https://github.com/abhishekkrthakur/diffuzers), and [lama-cleaner](https://github.com/Sanster/lama-cleaner)
- Every pipeline should have one and only one way to run it via a `__call__` method. The naming of the `__call__` arguments should be shared across all pipelines.
- Pipelines should be named after the task they are intended to solve.
- In almost all cases, novel diffusion pipelines shall be implemented in a new pipeline folder/file.
### Models
Models are designed as configurable toolboxes that are natural extensions of [PyTorch's Module class](https://pytorch.org/docs/stable/generated/torch.nn.Module.html). They only partly follow the **single-file policy**.
The following design principles are followed:
- Models correspond to **a type of model architecture**. *E.g.* the [`UNet2DConditionModel`] class is used for all UNet variations that expect 2D image inputs and are conditioned on some context.
- All models can be found in [`src/diffusers/models`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/models) and every model architecture shall be defined in its file, e.g. [`unet_2d_condition.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_condition.py), [`transformer_2d.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/transformer_2d.py), etc...
- Models **do not** follow the single-file policy and should make use of smaller model building blocks, such as [`attention.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/attention.py), [`resnet.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/resnet.py), [`embeddings.py`](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/embeddings.py), etc... **Note**: This is in stark contrast to Transformers' modeling files and shows that models do not really follow the single-file policy.
- Models intend to expose complexity, just like PyTorch's module does, and give clear error messages.
- Models all inherit from `ModelMixin` and `ConfigMixin`.
- Models can be optimized for performance when it doesnt demand major code changes, keeps backward compatibility, and gives significant memory or compute gain.
- Models should by default have the highest precision and lowest performance setting.
- To integrate new model checkpoints whose general architecture can be classified as an architecture that already exists in Diffusers, the existing model architecture shall be adapted to make it work with the new checkpoint. One should only create a new file if the model architecture is fundamentally different.
- Models should be designed to be easily extendable to future changes. This can be achieved by limiting public function arguments, configuration arguments, and "foreseeing" future changes, *e.g.* it is usually better to add `string` "...type" arguments that can easily be extended to new future types instead of boolean `is_..._type` arguments. Only the minimum amount of changes shall be made to existing architectures to make a new model checkpoint work.
- The model design is a difficult trade-off between keeping code readable and concise and supporting many model checkpoints. For most parts of the modeling code, classes shall be adapted for new model checkpoints, while there are some exceptions where it is preferred to add new classes to make sure the code is kept concise and
readable longterm, such as [UNet blocks](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/unet_2d_blocks.py) and [Attention processors](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/cross_attention.py).
### Schedulers
Schedulers are responsible to guide the denoising process for inference as well as to define a noise schedule for training. They are designed as individual classes with loadable configuration files and strongly follow the **single-file policy**.
The following design principles are followed:
- All schedulers are found in [`src/diffusers/schedulers`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/schedulers).
- Schedulers are **not** allowed to import from large utils files and shall be kept very self-contained.
- One scheduler python file corresponds to one scheduler algorithm (as might be defined in a paper).
- If schedulers share similar functionalities, we can make use of the `#Copied from` mechanism.
- Schedulers all inherit from `SchedulerMixin` and `ConfigMixin`.
- Schedulers can be easily swapped out with the [`ConfigMixin.from_config`](https://huggingface.co/docs/diffusers/main/en/api/configuration#diffusers.ConfigMixin.from_config) method as explained in detail [here](./using-diffusers/schedulers.mdx).
- Every scheduler has to have a `set_num_inference_steps`, and a `step` function. `set_num_inference_steps(...)` has to be called before every denoising process, *i.e.* before `step(...)` is called.
- Every scheduler exposes the timesteps to be "looped over" via a `timesteps` attribute, which is an array of timesteps the model will be called upon
- The `step(...)` function takes a predicted model output and the "current" sample (x_t) and returns the "previous", slightly more denoised sample (x_t-1).
- Given the complexity of diffusion schedulers, the `step` function does not expose all the complexity and can be a bit of a "black box".
- In almost all cases, novel schedulers shall be implemented in a new scheduling file.

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# How to use Stable Diffusion on Habana Gaudi
🤗 Diffusers is compatible with Habana Gaudi through 🤗 [Optimum Habana](https://huggingface.co/docs/optimum/habana/usage_guides/stable_diffusion).
## Requirements
- Optimum Habana 1.3 or later, [here](https://huggingface.co/docs/optimum/habana/installation) is how to install it.
- SynapseAI 1.7.
## Inference Pipeline
To generate images with Stable Diffusion 1 and 2 on Gaudi, you need to instantiate two instances:
- A pipeline with [`GaudiStableDiffusionPipeline`](https://huggingface.co/docs/optimum/habana/package_reference/stable_diffusion_pipeline). This pipeline supports *text-to-image generation*.
- A scheduler with [`GaudiDDIMScheduler`](https://huggingface.co/docs/optimum/habana/package_reference/stable_diffusion_pipeline#optimum.habana.diffusers.GaudiDDIMScheduler). This scheduler has been optimized for Habana Gaudi.
When initializing the pipeline, you have to specify `use_habana=True` to deploy it on HPUs.
Furthermore, in order to get the fastest possible generations you should enable **HPU graphs** with `use_hpu_graphs=True`.
Finally, you will need to specify a [Gaudi configuration](https://huggingface.co/docs/optimum/habana/package_reference/gaudi_config) which can be downloaded from the [Hugging Face Hub](https://huggingface.co/Habana).
```python
from optimum.habana import GaudiConfig
from optimum.habana.diffusers import GaudiDDIMScheduler, GaudiStableDiffusionPipeline
model_name = "stabilityai/stable-diffusion-2-base"
scheduler = GaudiDDIMScheduler.from_pretrained(model_name, subfolder="scheduler")
pipeline = GaudiStableDiffusionPipeline.from_pretrained(
model_name,
scheduler=scheduler,
use_habana=True,
use_hpu_graphs=True,
gaudi_config="Habana/stable-diffusion",
)
```
You can then call the pipeline to generate images by batches from one or several prompts:
```python
outputs = pipeline(
prompt=[
"High quality photo of an astronaut riding a horse in space",
"Face of a yellow cat, high resolution, sitting on a park bench",
],
num_images_per_prompt=10,
batch_size=4,
)
```
For more information, check out Optimum Habana's [documentation](https://huggingface.co/docs/optimum/habana/usage_guides/stable_diffusion) and the [example](https://github.com/huggingface/optimum-habana/tree/main/examples/stable-diffusion) provided in the official Github repository.
## Benchmark
Here are the latencies for Habana Gaudi 1 and Gaudi 2 with the [Habana/stable-diffusion](https://huggingface.co/Habana/stable-diffusion) Gaudi configuration (mixed precision bf16/fp32):
| | Latency | Batch size |
| ------- |:-------:|:----------:|
| Gaudi 1 | 4.37s | 4/8 |
| Gaudi 2 | 1.19s | 4/8 |

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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Installing xFormers
We recommend the use of [xFormers](https://github.com/facebookresearch/xformers) for both inference and training. In our tests, the optimizations performed in the attention blocks allow for both faster speed and reduced memory consumption.
Starting from version `0.0.16` of xFormers, released on January 2023, installation can be easily performed using pre-built pip wheels:
```bash
pip install xformers
```
<Tip>
The xFormers PIP package requires the latest version of PyTorch (1.13.1 as of xFormers 0.0.16). If you need to use a previous version of PyTorch, then we recommend you install xFormers from source using [the project instructions](https://github.com/facebookresearch/xformers#installing-xformers).
</Tip>
After xFormers is installed, you can use `enable_xformers_memory_efficient_attention()` for faster inference and reduced memory consumption, as discussed [here](fp16#memory-efficient-attention).
<Tip warning={true}>
According to [this issue](https://github.com/huggingface/diffusers/issues/2234#issuecomment-1416931212), xFormers `v0.0.16` cannot be used for training (fine-tune or Dreambooth) in some GPUs. If you observe that problem, please install a development version as indicated in that comment.
</Tip>

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
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the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# The Stable Diffusion Guide 🎨
<a target="_blank" href="https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/sd_101_guide.ipynb">
<img src="https://colab.research.google.com/assets/colab-badge.svg" alt="Open In Colab"/>
</a>
## Intro
Stable Diffusion is a [Latent Diffusion model](https://github.com/CompVis/latent-diffusion) developed by researchers from the Machine Vision and Learning group at LMU Munich, *a.k.a* CompVis.
Model checkpoints were publicly released at the end of August 2022 by a collaboration of Stability AI, CompVis, and Runway with support from EleutherAI and LAION. For more information, you can check out [the official blog post](https://stability.ai/blog/stable-diffusion-public-release).
Since its public release the community has done an incredible job at working together to make the stable diffusion checkpoints **faster**, **more memory efficient**, and **more performant**.
🧨 Diffusers offers a simple API to run stable diffusion with all memory, computing, and quality improvements.
This notebook walks you through the improvements one-by-one so you can best leverage [`StableDiffusionPipeline`] for **inference**.
## Prompt Engineering 🎨
When running *Stable Diffusion* in inference, we usually want to generate a certain type, or style of image and then improve upon it. Improving upon a previously generated image means running inference over and over again with a different prompt and potentially a different seed until we are happy with our generation.
So to begin with, it is most important to speed up stable diffusion as much as possible to generate as many pictures as possible in a given amount of time.
This can be done by both improving the **computational efficiency** (speed) and the **memory efficiency** (GPU RAM).
Let's start by looking into computational efficiency first.
Throughout the notebook, we will focus on [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5):
``` python
model_id = "runwayml/stable-diffusion-v1-5"
```
Let's load the pipeline.
## Speed Optimization
``` python
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(model_id)
```
We aim at generating a beautiful photograph of an *old warrior chief* and will later try to find the best prompt to generate such a photograph. For now, let's keep the prompt simple:
``` python
prompt = "portrait photo of a old warrior chief"
```
To begin with, we should make sure we run inference on GPU, so let's move the pipeline to GPU, just like you would with any PyTorch module.
``` python
pipe = pipe.to("cuda")
```
To generate an image, you should use the [~`StableDiffusionPipeline.__call__`] method.
To make sure we can reproduce more or less the same image in every call, let's make use of the generator. See the documentation on reproducibility [here](./conceptual/reproducibility) for more information.
``` python
generator = torch.Generator("cuda").manual_seed(0)
```
Now, let's take a spin on it.
``` python
image = pipe(prompt, generator=generator).images[0]
image
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_1.png)
Cool, this now took roughly 30 seconds on a T4 GPU (you might see faster inference if your allocated GPU is better than a T4).
The default run we did above used full float32 precision and ran the default number of inference steps (50). The easiest speed-ups come from switching to float16 (or half) precision and simply running fewer inference steps. Let's load the model now in float16 instead.
``` python
import torch
pipe = StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
pipe = pipe.to("cuda")
```
And we can again call the pipeline to generate an image.
``` python
generator = torch.Generator("cuda").manual_seed(0)
image = pipe(prompt, generator=generator).images[0]
image
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_2.png)
Cool, this is almost three times as fast for arguably the same image quality.
We strongly suggest always running your pipelines in float16 as so far we have very rarely seen degradations in quality because of it.
Next, let's see if we need to use 50 inference steps or whether we could use significantly fewer. The number of inference steps is associated with the denoising scheduler we use. Choosing a more efficient scheduler could help us decrease the number of steps.
Let's have a look at all the schedulers the stable diffusion pipeline is compatible with.
``` python
pipe.scheduler.compatibles
```
```
[diffusers.schedulers.scheduling_dpmsolver_singlestep.DPMSolverSinglestepScheduler,
diffusers.schedulers.scheduling_lms_discrete.LMSDiscreteScheduler,
diffusers.schedulers.scheduling_heun_discrete.HeunDiscreteScheduler,
diffusers.schedulers.scheduling_pndm.PNDMScheduler,
diffusers.schedulers.scheduling_euler_discrete.EulerDiscreteScheduler,
diffusers.schedulers.scheduling_euler_ancestral_discrete.EulerAncestralDiscreteScheduler,
diffusers.schedulers.scheduling_dpmsolver_multistep.DPMSolverMultistepScheduler,
diffusers.schedulers.scheduling_ddpm.DDPMScheduler,
diffusers.schedulers.scheduling_ddim.DDIMScheduler]
```
Cool, that's a lot of schedulers.
🧨 Diffusers is constantly adding a bunch of novel schedulers/samplers that can be used with Stable Diffusion. For more information, we recommend taking a look at the official documentation [here](https://huggingface.co/docs/diffusers/main/en/api/schedulers/overview).
Alright, right now Stable Diffusion is using the `PNDMScheduler` which usually requires around 50 inference steps. However, other schedulers such as `DPMSolverMultistepScheduler` or `DPMSolverSinglestepScheduler` seem to get away with just 20 to 25 inference steps. Let's try them out.
You can set a new scheduler by making use of the [from_config](https://huggingface.co/docs/diffusers/main/en/api/configuration#diffusers.ConfigMixin.from_config) function.
``` python
from diffusers import DPMSolverMultistepScheduler
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
```
Now, let's try to reduce the number of inference steps to just 20.
``` python
generator = torch.Generator("cuda").manual_seed(0)
image = pipe(prompt, generator=generator, num_inference_steps=20).images[0]
image
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_3.png)
The image now does look a little different, but it's arguably still of equally high quality. We now cut inference time to just 4 seconds though 😍.
## Memory Optimization
Less memory used in generation indirectly implies more speed, since we're often trying to maximize how many images we can generate per second. Usually, the more images per inference run, the more images per second too.
The easiest way to see how many images we can generate at once is to simply try it out, and see when we get a *"Out-of-memory (OOM)"* error.
We can run batched inference by simply passing a list of prompts and generators. Let's define a quick function that generates a batch for us.
``` python
def get_inputs(batch_size=1):
generator = [torch.Generator("cuda").manual_seed(i) for i in range(batch_size)]
prompts = batch_size * [prompt]
num_inference_steps = 20
return {"prompt": prompts, "generator": generator, "num_inference_steps": num_inference_steps}
```
This function returns a list of prompts and a list of generators, so we can reuse the generator that produced a result we like.
We also need a method that allows us to easily display a batch of images.
``` python
from PIL import Image
def image_grid(imgs, rows=2, cols=2):
w, h = imgs[0].size
grid = Image.new('RGB', size=(cols*w, rows*h))
for i, img in enumerate(imgs):
grid.paste(img, box=(i%cols*w, i//cols*h))
return grid
```
Cool, let's see how much memory we can use starting with `batch_size=4`.
``` python
images = pipe(**get_inputs(batch_size=4)).images
image_grid(images)
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_4.png)
Going over a batch_size of 4 will error out in this notebook (assuming we are running it on a T4 GPU). Also, we can see we only generate slightly more images per second (3.75s/image) compared to 4s/image previously.
However, the community has found some nice tricks to improve the memory constraints further. After stable diffusion was released, the community found improvements within days and shared them freely over GitHub - open-source at its finest! I believe the original idea came from [this](https://github.com/basujindal/stable-diffusion/pull/117) GitHub thread.
By far most of the memory is taken up by the cross-attention layers. Instead of running this operation in batch, one can run it sequentially to save a significant amount of memory.
It can easily be enabled by calling `enable_attention_slicing` as is documented [here](https://huggingface.co/docs/diffusers/main/en/api/pipelines/stable_diffusion/text2img#diffusers.StableDiffusionPipeline.enable_attention_slicing).
``` python
pipe.enable_attention_slicing()
```
Great, now that attention slicing is enabled, let's try to double the batch size again, going for `batch_size=8`.
``` python
images = pipe(**get_inputs(batch_size=8)).images
image_grid(images, rows=2, cols=4)
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_5.png)
Nice, it works. However, the speed gain is again not very big (it might however be much more significant on other GPUs).
We're at roughly 3.5 seconds per image 🔥 which is probably the fastest we can be with a simple T4 without sacrificing quality.
Next, let's look into how to improve the quality!
## Quality Improvements
Now that our image generation pipeline is blazing fast, let's try to get maximum image quality.
First of all, image quality is extremely subjective, so it's difficult to make general claims here.
The most obvious step to take to improve quality is to use *better checkpoints*. Since the release of Stable Diffusion, many improved versions have been released, which are summarized here:
- *Official Release - 22 Aug 2022*: [Stable-Diffusion 1.4](https://huggingface.co/CompVis/stable-diffusion-v1-4)
- *20 October 2022*: [Stable-Diffusion 1.5](https://huggingface.co/runwayml/stable-diffusion-v1-5)
- *24 Nov 2022*: [Stable-Diffusion 2.0](https://huggingface.co/stabilityai/stable-diffusion-2-0)
- *7 Dec 2022*: [Stable-Diffusion 2.1](https://huggingface.co/stabilityai/stable-diffusion-2-1)
Newer versions don't necessarily mean better image quality with the same parameters. People mentioned that *2.0* is slightly worse than *1.5* for certain prompts, but given the right prompt engineering *2.0* and *2.1* seem to be better.
Overall, we strongly recommend just trying the models out and reading up on advice online (e.g. it has been shown that using negative prompts is very important for 2.0 and 2.1 to get the highest possible quality. See for example [this nice blog post](https://minimaxir.com/2022/11/stable-diffusion-negative-prompt/).
Additionally, the community has started fine-tuning many of the above versions on certain styles with some of them having an extremely high quality and gaining a lot of traction.
We recommend having a look at all [diffusers checkpoints sorted by downloads and trying out the different checkpoints](https://huggingface.co/models?library=diffusers).
For the following, we will stick to v1.5 for simplicity.
Next, we can also try to optimize single components of the pipeline, e.g. switching out the latent decoder. For more details on how the whole Stable Diffusion pipeline works, please have a look at [this blog post](https://huggingface.co/blog/stable_diffusion).
Let's load [stabilityai's newest auto-decoder](https://huggingface.co/stabilityai/stable-diffusion-2-1).
``` python
from diffusers import AutoencoderKL
vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse", torch_dtype=torch.float16).to("cuda")
```
Now we can set it to the vae of the pipeline to use it.
``` python
pipe.vae = vae
```
Let's run the same prompt as before to compare quality.
``` python
images = pipe(**get_inputs(batch_size=8)).images
image_grid(images, rows=2, cols=4)
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_6.png)
Seems like the difference is only very minor, but the new generations are arguably a bit *sharper*.
Cool, finally, let's look a bit into prompt engineering.
Our goal was to generate a photo of an old warrior chief. Let's now try to bring a bit more color into the photos and make the look more impressive.
Originally our prompt was "*portrait photo of an old warrior chief*".
To improve the prompt, it often helps to add cues that could have been used online to save high-quality photos, as well as add more details.
Essentially, when doing prompt engineering, one has to think:
- How was the photo or similar photos of the one I want probably stored on the internet?
- What additional detail can I give that steers the models into the style that I want?
Cool, let's add more details.
``` python
prompt += ", tribal panther make up, blue on red, side profile, looking away, serious eyes"
```
and let's also add some cues that usually help to generate higher quality images.
``` python
prompt += " 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta"
prompt
```
Cool, let's now try this prompt.
``` python
images = pipe(**get_inputs(batch_size=8)).images
image_grid(images, rows=2, cols=4)
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_7.png)
Pretty impressive! We got some very high-quality image generations there. The 2nd image is my personal favorite, so I'll re-use this seed and see whether I can tweak the prompts slightly by using "oldest warrior", "old", "", and "young" instead of "old".
``` python
prompts = [
"portrait photo of the oldest warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
"portrait photo of a old warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
"portrait photo of a warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
"portrait photo of a young warrior chief, tribal panther make up, blue on red, side profile, looking away, serious eyes 50mm portrait photography, hard rim lighting photography--beta --ar 2:3 --beta --upbeta",
]
generator = [torch.Generator("cuda").manual_seed(1) for _ in range(len(prompts))] # 1 because we want the 2nd image
images = pipe(prompt=prompts, generator=generator, num_inference_steps=25).images
image_grid(images)
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/stable_diffusion_101/sd_101_8.png)
The first picture looks nice! The eye movement slightly changed and looks nice. This finished up our 101-guide on how to use Stable Diffusion 🤗.
For more information on optimization or other guides, I recommend taking a look at the following:
- [Blog post about Stable Diffusion](https://huggingface.co/blog/stable_diffusion): In-detail blog post explaining Stable Diffusion.
- [FlashAttention](https://huggingface.co/docs/diffusers/optimization/xformers): XFormers flash attention can optimize your model even further with more speed and memory improvements.
- [Dreambooth](https://huggingface.co/docs/diffusers/training/dreambooth) - Quickly customize the model by fine-tuning it.
- [General info on Stable Diffusion](https://huggingface.co/docs/diffusers/main/en/api/pipelines/stable_diffusion/overview) - Info on other tasks that are powered by Stable Diffusion.

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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# LoRA Support in Diffusers
Diffusers supports LoRA for faster fine-tuning of Stable Diffusion, allowing greater memory efficiency and easier portability.
Low-Rank Adaption of Large Language Models was first introduced by Microsoft in
[LoRA: Low-Rank Adaptation of Large Language Models](https://arxiv.org/abs/2106.09685) by *Edward J. Hu, Yelong Shen, Phillip Wallis, Zeyuan Allen-Zhu, Yuanzhi Li, Shean Wang, Lu Wang, Weizhu Chen*.
In a nutshell, LoRA allows adapting pretrained models by adding pairs of rank-decomposition weight matrices (called **update matrices**)
to existing weights and **only** training those newly added weights. This has a couple of advantages:
- Previous pretrained weights are kept frozen so that the model is not so prone to [catastrophic forgetting](https://www.pnas.org/doi/10.1073/pnas.1611835114).
- Rank-decomposition matrices have significantly fewer parameters than the original model, which means that trained LoRA weights are easily portable.
- LoRA matrices are generally added to the attention layers of the original model and they control to which extent the model is adapted toward new training images via a `scale` parameter.
**__Note that the usage of LoRA is not just limited to attention layers. In the original LoRA work, the authors found out that just amending
the attention layers of a language model is sufficient to obtain good downstream performance with great efficiency. This is why, it's common
to just add the LoRA weights to the attention layers of a model.__**
[cloneofsimo](https://github.com/cloneofsimo) was the first to try out LoRA training for Stable Diffusion in the popular [lora](https://github.com/cloneofsimo/lora) GitHub repository.
<Tip>
LoRA allows us to achieve greater memory efficiency since the pretrained weights are kept frozen and only the LoRA weights are trained, thereby
allowing us to run fine-tuning on consumer GPUs like Tesla T4, RTX 3080 or even RTX 2080 Ti! One can get access to GPUs like T4 in the free
tiers of Kaggle Kernels and Google Colab Notebooks.
</Tip>
## Getting started with LoRA for fine-tuning
Stable Diffusion can be fine-tuned in different ways:
* [Textual inversion](https://huggingface.co/docs/diffusers/main/en/training/text_inversion)
* [DreamBooth](https://huggingface.co/docs/diffusers/main/en/training/dreambooth)
* [Text2Image fine-tuning](https://huggingface.co/docs/diffusers/main/en/training/text2image)
We provide two end-to-end examples that show how to run fine-tuning with LoRA:
* [DreamBooth](https://github.com/huggingface/diffusers/tree/main/examples/dreambooth#training-with-low-rank-adaptation-of-large-language-models-lora)
* [Text2Image](https://github.com/huggingface/diffusers/tree/main/examples/text_to_image#training-with-lora)
If you want to perform DreamBooth training with LoRA, for instance, you would run:
```bash
export MODEL_NAME="runwayml/stable-diffusion-v1-5"
export INSTANCE_DIR="path-to-instance-images"
export OUTPUT_DIR="path-to-save-model"
accelerate launch train_dreambooth_lora.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--output_dir=$OUTPUT_DIR \
--instance_prompt="a photo of sks dog" \
--resolution=512 \
--train_batch_size=1 \
--gradient_accumulation_steps=1 \
--checkpointing_steps=100 \
--learning_rate=1e-4 \
--report_to="wandb" \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--max_train_steps=500 \
--validation_prompt="A photo of sks dog in a bucket" \
--validation_epochs=50 \
--seed="0" \
--push_to_hub
```
A similar process can be followed to fully fine-tune Stable Diffusion on a custom dataset using the
`examples/text_to_image/train_text_to_image_lora.py` script.
Refer to the respective examples linked above to learn more.
<Tip>
When using LoRA we can use a much higher learning rate (typically 1e-4 as opposed to ~1e-6) compared to non-LoRA Dreambooth fine-tuning.
</Tip>
But there is no free lunch. For the given dataset and expected generation quality, you'd still need to experiment with
different hyperparameters. Here are some important ones:
* Training time
* Learning rate
* Number of training steps
* Inference time
* Number of steps
* Scheduler type
Additionally, you can follow [this blog](https://huggingface.co/blog/dreambooth) that documents some of our experimental
findings for performing DreamBooth training of Stable Diffusion.
When fine-tuning, the LoRA update matrices are only added to the attention layers. To enable this, we added new weight
loading functionalities. Their details are available [here](https://huggingface.co/docs/diffusers/main/en/api/loaders).
## Inference
Assuming you used the `examples/text_to_image/train_text_to_image_lora.py` to fine-tune Stable Diffusion on the [Pokemon
dataset](https://huggingface.co/datasets/lambdalabs/pokemon-blip-captions), you can perform inference like so:
```py
from diffusers import StableDiffusionPipeline
import torch
model_path = "sayakpaul/sd-model-finetuned-lora-t4"
pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16)
pipe.unet.load_attn_procs(model_path)
pipe.to("cuda")
prompt = "A pokemon with blue eyes."
image = pipe(prompt, num_inference_steps=30, guidance_scale=7.5).images[0]
image.save("pokemon.png")
```
Here are some example images you can expect:
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/pokemon-collage.png"/>
[`sayakpaul/sd-model-finetuned-lora-t4`](https://huggingface.co/sayakpaul/sd-model-finetuned-lora-t4) contains [LoRA fine-tuned update matrices](https://huggingface.co/sayakpaul/sd-model-finetuned-lora-t4/blob/main/pytorch_lora_weights.bin)
which is only 3 MBs in size. During inference, the pre-trained Stable Diffusion checkpoints are loaded alongside these update
matrices and then they are combined to run inference.
You can use the [`huggingface_hub`](https://github.com/huggingface/huggingface_hub) library to retrieve the base model
from [`sayakpaul/sd-model-finetuned-lora-t4`](https://huggingface.co/sayakpaul/sd-model-finetuned-lora-t4) like so:
```py
from huggingface_hub.repocard import RepoCard
card = RepoCard.load("sayakpaul/sd-model-finetuned-lora-t4")
base_model = card.data.to_dict()["base_model"]
# 'CompVis/stable-diffusion-v1-4'
```
And then you can use `pipe = StableDiffusionPipeline.from_pretrained(base_model, torch_dtype=torch.float16)`.
This is especially useful when you don't want to hardcode the base model identifier during initializing the `StableDiffusionPipeline`.
Inference for DreamBooth training remains the same. Check
[this section](https://github.com/huggingface/diffusers/tree/main/examples/dreambooth#inference-1) for more details.
### Merging LoRA with original model
When performing inference, you can merge the trained LoRA weights with the frozen pre-trained model weights, to interpolate between the original model's inference result (as if no fine-tuning had occurred) and the fully fine-tuned version.
You can adjust the merging ratio with a parameter called α (alpha) in the paper, or `scale` in our implementation. You can tweak it with the following code, that passes `scale` as `cross_attention_kwargs` in the pipeline call:
```py
from diffusers import StableDiffusionPipeline
import torch
model_path = "sayakpaul/sd-model-finetuned-lora-t4"
pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16)
pipe.unet.load_attn_procs(model_path)
pipe.to("cuda")
prompt = "A pokemon with blue eyes."
image = pipe(prompt, num_inference_steps=30, guidance_scale=7.5, cross_attention_kwargs={"scale": 0.5}).images[0]
image.save("pokemon.png")
```
A value of `0` is the same as _not_ using the LoRA weights, whereas `1` means only the LoRA fine-tuned weights will be used. Values between 0 and 1 will interpolate between the two versions.
## Known limitations
* Currently, we only support LoRA for the attention layers of [`UNet2DConditionModel`](https://huggingface.co/docs/diffusers/main/en/api/models#diffusers.UNet2DConditionModel).

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@@ -1,35 +0,0 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Text-Guided Image-to-Image Generation
The [`StableDiffusionDepth2ImgPipeline`] lets you pass a text prompt and an initial image to condition the generation of new images as well as a `depth_map` to preserve the images' structure. If no `depth_map` is provided, the pipeline will automatically predict the depth via an integrated depth-estimation model.
```python
import torch
import requests
from PIL import Image
from diffusers import StableDiffusionDepth2ImgPipeline
pipe = StableDiffusionDepth2ImgPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-depth",
torch_dtype=torch.float16,
).to("cuda")
url = "http://images.cocodataset.org/val2017/000000039769.jpg"
init_image = Image.open(requests.get(url, stream=True).raw)
prompt = "two tigers"
n_prompt = "bad, deformed, ugly, bad anatomy"
image = pipe(prompt=prompt, image=init_image, negative_prompt=n_prompt, strength=0.7).images[0]
```

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<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Using KerasCV Stable Diffusion Checkpoints in Diffusers
<Tip warning={true}>
This is an experimental feature.
</Tip>
[KerasCV](https://github.com/keras-team/keras-cv/) provides APIs for implementing various computer vision workflows. It
also provides the Stable Diffusion [v1 and v2](https://github.com/keras-team/keras-cv/blob/master/keras_cv/models/stable_diffusion)
models. Many practitioners find it easy to fine-tune the Stable Diffusion models shipped by KerasCV. However, as of this writing, KerasCV offers limited support to experiment with Stable Diffusion models for inference and deployment. On the other hand,
Diffusers provides tooling dedicated to this purpose (and more), such as different [noise schedulers](https://huggingface.co/docs/diffusers/using-diffusers/schedulers), [flash attention](https://huggingface.co/docs/diffusers/optimization/xformers), and [other
optimization techniques](https://huggingface.co/docs/diffusers/optimization/fp16).
How about fine-tuning Stable Diffusion models in KerasCV and exporting them such that they become compatible with Diffusers to combine the
best of both worlds? We have created a [tool](https://huggingface.co/spaces/sayakpaul/convert-kerascv-sd-diffusers) that
lets you do just that! It takes KerasCV Stable Diffusion checkpoints and exports them to Diffusers-compatible checkpoints.
More specifically, it first converts the checkpoints to PyTorch and then wraps them into a
[`StableDiffusionPipeline`](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/overview) which is ready
for inference. Finally, it pushes the converted checkpoints to a repository on the Hugging Face Hub.
We welcome you to try out the tool [here](https://huggingface.co/spaces/sayakpaul/convert-kerascv-sd-diffusers)
and share feedback via [discussions](https://huggingface.co/spaces/sayakpaul/convert-kerascv-sd-diffusers/discussions/new).
## Getting Started
First, you need to obtain the fine-tuned KerasCV Stable Diffusion checkpoints. We provide an
overview of the different ways Stable Diffusion models can be fine-tuned [using `diffusers`](https://huggingface.co/docs/diffusers/training/overview). For the Keras implementation of some of these methods, you can check out these resources:
* [Teach StableDiffusion new concepts via Textual Inversion](https://keras.io/examples/generative/fine_tune_via_textual_inversion/)
* [Fine-tuning Stable Diffusion](https://keras.io/examples/generative/finetune_stable_diffusion/)
* [DreamBooth](https://keras.io/examples/generative/dreambooth/)
* [Prompt-to-Prompt editing](https://github.com/miguelCalado/prompt-to-prompt-tensorflow)
Stable Diffusion is comprised of the following models:
* Text encoder
* UNet
* VAE
Depending on the fine-tuning task, we may fine-tune one or more of these components (the VAE is almost always left untouched). Here are some common combinations:
* DreamBooth: UNet and text encoder
* Classical text to image fine-tuning: UNet
* Textual Inversion: Just the newly initialized embeddings in the text encoder
### Performing the Conversion
Let's use [this checkpoint](https://huggingface.co/sayakpaul/textual-inversion-kerasio/resolve/main/textual_inversion_kerasio.h5) which was generated
by conducting Textual Inversion with the following "placeholder token": `<my-funny-cat-token>`.
On the tool, we supply the following things:
* Path(s) to download the fine-tuned checkpoint(s) (KerasCV)
* An HF token
* Placeholder token (only applicable for Textual Inversion)
<div align="center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/space_snap.png"/>
</div>
As soon as you hit "Submit", the conversion process will begin. Once it's complete, you should see the following:
<div align="center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/model_push_success.png"/>
</div>
If you click the [link](https://huggingface.co/sayakpaul/textual-inversion-cat-kerascv_sd_diffusers_pipeline/tree/main), you
should see something like so:
<div align="center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/model_repo_contents.png"/>
</div>
If you head over to the [model card of the repository](https://huggingface.co/sayakpaul/textual-inversion-cat-kerascv_sd_diffusers_pipeline), the
following should appear:
<div align="center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/model_card.png"/>
</div>
<Tip>
Note that we're not specifying the UNet weights here since the UNet is not fine-tuned during Textual Inversion.
</Tip>
And that's it! You now have your fine-tuned KerasCV Stable Diffusion model in Diffusers 🧨
## Using the Converted Model in Diffusers
Just beside the model card of the [repository](https://huggingface.co/sayakpaul/textual-inversion-cat-kerascv_sd_diffusers_pipeline),
you'd notice an inference widget to try out the model directly from the UI 🤗
<div align="center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/inference_widget_output.png"/>
</div>
On the top right hand side, we provide a "Use in Diffusers" button. If you click the button, you should see the following code-snippet:
```py
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained("sayakpaul/textual-inversion-cat-kerascv_sd_diffusers_pipeline")
```
The model is in standard `diffusers` format. Let's perform inference!
```py
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained("sayakpaul/textual-inversion-cat-kerascv_sd_diffusers_pipeline")
pipeline.to("cuda")
placeholder_token = "<my-funny-cat-token>"
prompt = f"two {placeholder_token} getting married, photorealistic, high quality"
image = pipeline(prompt, num_inference_steps=50).images[0]
```
And we get:
<div align="center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/diffusers_output_one.png"/>
</div>
_**Note that if you specified a `placeholder_token` while performing the conversion, the tool will log it accordingly. Refer
to the model card of [this repository](https://huggingface.co/sayakpaul/textual-inversion-cat-kerascv_sd_diffusers_pipeline)
as an example.**_
We welcome you to use the tool for various Stable Diffusion fine-tuning scenarios and let us know your feedback! Here are some examples
of Diffusers checkpoints that were obtained using the tool:
* [sayakpaul/text-unet-dogs-kerascv_sd_diffusers_pipeline](https://huggingface.co/sayakpaul/text-unet-dogs-kerascv_sd_diffusers_pipeline) (DreamBooth with both the text encoder and UNet fine-tuned)
* [sayakpaul/unet-dogs-kerascv_sd_diffusers_pipeline](https://huggingface.co/sayakpaul/unet-dogs-kerascv_sd_diffusers_pipeline) (DreamBooth with only the UNet fine-tuned)
## Incorporating Diffusers Goodies 🎁
Diffusers provides various options that one can leverage to experiment with different inference setups. One particularly
useful option is the use of a different noise scheduler during inference other than what was used during fine-tuning.
Let's try out the [`DPMSolverMultistepScheduler`](https://huggingface.co/docs/diffusers/main/en/api/schedulers/multistep_dpm_solver)
which is different from the one ([`DDPMScheduler`](https://huggingface.co/docs/diffusers/main/en/api/schedulers/ddpm)) used during
fine-tuning.
You can read more details about this process in [this section](https://huggingface.co/docs/diffusers/using-diffusers/schedulers).
```py
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler
pipeline = DiffusionPipeline.from_pretrained("sayakpaul/textual-inversion-cat-kerascv_sd_diffusers_pipeline")
pipeline.scheduler = DPMSolverMultistepScheduler.from_config(pipeline.scheduler.config)
pipeline.to("cuda")
placeholder_token = "<my-funny-cat-token>"
prompt = f"two {placeholder_token} getting married, photorealistic, high quality"
image = pipeline(prompt, num_inference_steps=50).images[0]
```
<div align="center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/diffusers_output_two.png"/>
</div>
One can also continue fine-tuning from these Diffusers checkpoints by leveraging some relevant tools from Diffusers. Refer [here](https://huggingface.co/docs/diffusers/training/overview) for
more details. For inference-specific optimizations, refer [here](https://huggingface.co/docs/diffusers/main/en/optimization/fp16).
## Known Limitations
* Only Stable Diffusion v1 checkpoints are supported for conversion in this tool.

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Reproducibility
Before reading about reproducibility for Diffusers, it is strongly recommended to take a look at
[PyTorch's statement about reproducibility](https://pytorch.org/docs/stable/notes/randomness.html).
PyTorch states that
> *completely reproducible results are not guaranteed across PyTorch releases, individual commits, or different platforms.*
While one can never expect the same results across platforms, one can expect results to be reproducible
across releases, platforms, etc... within a certain tolerance. However, this tolerance strongly varies
depending on the diffusion pipeline and checkpoint.
In the following, we show how to best control sources of randomness for diffusion models.
## Inference
During inference, diffusion pipelines heavily rely on random sampling operations, such as the creating the
gaussian noise tensors to be denoised and adding noise to the scheduling step.
Let's have a look at an example. We run the [DDIM pipeline](./api/pipelines/ddim.mdx)
for just two inference steps and return a numpy tensor to look into the numerical values of the output.
```python
from diffusers import DDIMPipeline
import numpy as np
model_id = "google/ddpm-cifar10-32"
# load model and scheduler
ddim = DDIMPipeline.from_pretrained(model_id)
# run pipeline for just two steps and return numpy tensor
image = ddim(num_inference_steps=2, output_type="np").images
print(np.abs(image).sum())
```
Running the above prints a value of 1464.2076, but running it again prints a different
value of 1495.1768. What is going on here? Every time the pipeline is run, gaussian noise
is created and step-wise denoised. To create the gaussian noise with [`torch.randn`](https://pytorch.org/docs/stable/generated/torch.randn.html), a different random seed is taken every time, thus leading to a different result.
This is a desired property of diffusion pipelines, as it means that the pipeline can create a different random image every time it is run. In many cases, one would like to generate the exact same image of a certain
run, for which case an instance of a [PyTorch generator](https://pytorch.org/docs/stable/generated/torch.randn.html) has to be passed:
```python
import torch
from diffusers import DDIMPipeline
import numpy as np
model_id = "google/ddpm-cifar10-32"
# load model and scheduler
ddim = DDIMPipeline.from_pretrained(model_id)
# create a generator for reproducibility
generator = torch.Generator(device="cpu").manual_seed(0)
# run pipeline for just two steps and return numpy tensor
image = ddim(num_inference_steps=2, output_type="np", generator=generator).images
print(np.abs(image).sum())
```
Running the above always prints a value of 1491.1711 - also upon running it again because we
define the generator object to be passed to all random functions of the pipeline.
If you run this code snippet on your specific hardware and version, you should get a similar, if not the same, result.
<Tip>
It might be a bit unintuitive at first to pass `generator` objects to the pipelines instead of
just integer values representing the seed, but this is the recommended design when dealing with
probabilistic models in PyTorch as generators are *random states* that are advanced and can thus be
passed to multiple pipelines in a sequence.
</Tip>
Great! Now, we know how to write reproducible pipelines, but it gets a bit trickier since the above example only runs on the CPU. How do we also achieve reproducibility on GPU?
In short, one should not expect full reproducibility across different hardware when running pipelines on GPU
as matrix multiplications are less deterministic on GPU than on CPU and diffusion pipelines tend to require
a lot of matrix multiplications. Let's see what we can do to keep the randomness within limits across
different GPU hardware.
To achieve maximum speed performance, it is recommended to create the generator directly on GPU when running
the pipeline on GPU:
```python
import torch
from diffusers import DDIMPipeline
import numpy as np
model_id = "google/ddpm-cifar10-32"
# load model and scheduler
ddim = DDIMPipeline.from_pretrained(model_id)
ddim.to("cuda")
# create a generator for reproducibility
generator = torch.Generator(device="cuda").manual_seed(0)
# run pipeline for just two steps and return numpy tensor
image = ddim(num_inference_steps=2, output_type="np", generator=generator).images
print(np.abs(image).sum())
```
Running the above now prints a value of 1389.8634 - even though we're using the exact same seed!
This is unfortunate as it means we cannot reproduce the results we achieved on GPU, also on CPU.
Nevertheless, it should be expected since the GPU uses a different random number generator than the CPU.
To circumvent this problem, we created a [`randn_tensor`](#diffusers.utils.randn_tensor) function, which can create random noise
on the CPU and then move the tensor to GPU if necessary. The function is used everywhere inside the pipelines allowing the user to **always** pass a CPU generator even if the pipeline is run on GPU:
```python
import torch
from diffusers import DDIMPipeline
import numpy as np
model_id = "google/ddpm-cifar10-32"
# load model and scheduler
ddim = DDIMPipeline.from_pretrained(model_id)
ddim.to("cuda")
# create a generator for reproducibility
generator = torch.manual_seed(0)
# run pipeline for just two steps and return numpy tensor
image = ddim(num_inference_steps=2, output_type="np", generator=generator).images
print(np.abs(image).sum())
```
Running the above now prints a value of 1491.1713, much closer to the value of 1491.1711 when
the pipeline is fully run on the CPU.
<Tip>
As a consequence, we recommend always passing a CPU generator if Reproducibility is important.
The loss of performance is often neglectable, but one can be sure to generate much more similar
values than if the pipeline would have been run on CPU.
</Tip>
Finally, we noticed that more complex pipelines, such as [`UnCLIPPipeline`] are often extremely
susceptible to precision error propagation and thus one cannot expect even similar results across
different GPU hardware or PyTorch versions. In such cases, one has to make sure to run
exactly the same hardware and PyTorch version for full Reproducibility.
## Randomness utilities
### randn_tensor
[[autodoc]] diffusers.utils.randn_tensor

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Re-using seeds for fast prompt engineering
A common use case when generating images is to generate a batch of images, select one image and improve it with a better, more detailed prompt in a second run.
To do this, one needs to make each generated image of the batch deterministic.
Images are generated by denoising gaussian random noise which can be instantiated by passing a [torch generator](https://pytorch.org/docs/stable/generated/torch.Generator.html#generator).
Now, for batched generation, we need to make sure that every single generated image in the batch is tied exactly to one seed. In 🧨 Diffusers, this can be achieved by not passing one `generator`, but a list
of `generators` to the pipeline.
Let's go through an example using [`runwayml/stable-diffusion-v1-5`](runwayml/stable-diffusion-v1-5).
We want to generate several versions of the prompt:
```py
prompt = "Labrador in the style of Vermeer"
```
Let's load the pipeline
```python
>>> from diffusers import DiffusionPipeline
>>> pipe = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = pipe.to("cuda")
```
Now, let's define 4 different generators, since we would like to reproduce a certain image. We'll use seeds `0` to `3` to create our generators.
```python
>>> import torch
>>> generator = [torch.Generator(device="cuda").manual_seed(i) for i in range(4)]
```
Let's generate 4 images:
```python
>>> images = pipe(prompt, generator=generator, num_images_per_prompt=4).images
>>> images
```
![img](https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/reusabe_seeds.jpg)
Ok, the last images has some double eyes, but the first image looks good!
Let's try to make the prompt a bit better **while keeping the first seed**
so that the images are similar to the first image.
```python
prompt = [prompt + t for t in [", highly realistic", ", artsy", ", trending", ", colorful"]]
generator = [torch.Generator(device="cuda").manual_seed(0) for i in range(4)]
```
We create 4 generators with seed `0`, which is the first seed we used before.
Let's run the pipeline again.
```python
>>> images = pipe(prompt, generator=generator).images
>>> images
```
![img](https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/reusabe_seeds_2.jpg)

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@@ -1,19 +0,0 @@
# What is safetensors ?
[safetensors](https://github.com/huggingface/safetensors) is a different format
from the classic `.bin` which uses Pytorch which uses pickle.
Pickle is notoriously unsafe which allow any malicious file to execute arbitrary code.
The hub itself tries to prevent issues from it, but it's not a silver bullet.
`safetensors` first and foremost goal is to make loading machine learning models *safe*
in the sense that no takeover of your computer can be done.
# Why use safetensors ?
**Safety** can be one reason, if you're attempting to use a not well known model and
you're not sure about the source of the file.
And a secondary reason, is **the speed of loading**. Safetensors can load models much faster
than regular pickle files. If you spend a lot of times switching models, this can be
a huge timesave.

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@@ -1,87 +0,0 @@
# What is safetensors ?
[safetensors](https://github.com/huggingface/safetensors) is a different format
from the classic `.bin` which uses Pytorch which uses pickle. It contains the
exact same data, which is just the model weights (or tensors).
Pickle is notoriously unsafe which allow any malicious file to execute arbitrary code.
The hub itself tries to prevent issues from it, but it's not a silver bullet.
`safetensors` first and foremost goal is to make loading machine learning models *safe*
in the sense that no takeover of your computer can be done.
Hence the name.
# Why use safetensors ?
**Safety** can be one reason, if you're attempting to use a not well known model and
you're not sure about the source of the file.
And a secondary reason, is **the speed of loading**. Safetensors can load models much faster
than regular pickle files. If you spend a lot of times switching models, this can be
a huge timesave.
Numbers taken AMD EPYC 7742 64-Core Processor
```
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1")
# Loaded in safetensors 0:00:02.033658
# Loaded in Pytorch 0:00:02.663379
```
This is for the entire loading time, the actual weights loading time to load 500MB:
```
Safetensors: 3.4873ms
PyTorch: 172.7537ms
```
Performance in general is a tricky business, and there are a few things to understand:
- If you're using the model for the first time from the hub, you will have to download the weights.
That's extremely likely to be much slower than any loading method, therefore you will not see any difference
- If you're loading the model for the first time (let's say after a reboot) then your machine will have to
actually read the disk. It's likely to be as slow in both cases. Again the speed difference may not be as visible (this depends on hardware and the actual model).
- The best performance benefit is when the model was already loaded previously on your computer and you're switching from one model to another. Your OS, is trying really hard not to read from disk, since this is slow, so it will keep the files around in RAM, making it loading again much faster. Since safetensors is doing zero-copy of the tensors, reloading will be faster than pytorch since it has at least once extra copy to do.
# How to use safetensors ?
If you have `safetensors` installed, and all the weights are available in `safetensors` format, \
then by default it will use that instead of the pytorch weights.
If you are really paranoid about this, the ultimate weapon would be disabling `torch.load`:
```python
import torch
def _raise():
raise RuntimeError("I don't want to use pickle")
torch.load = lambda *args, **kwargs: _raise()
```
# I want to use model X but it doesn't have safetensors weights.
Just go to this [space](https://huggingface.co/spaces/safetensors/convert).
This will create a new PR with the weights, let's say `refs/pr/22`.
This space will download the pickled version, convert it, and upload it on the hub as a PR.
If anything bad is contained in the file, it's Huggingface hub that will get issues, not your own computer.
And we're equipped with dealing with it.
Then in order to use the model, even before the branch gets accepted by the original author you can do:
```python
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1", revision="refs/pr/22")
```
And that's it !
Anything unclear, concerns, or found a bugs ? [Open an issue](https://github.com/huggingface/diffusers/issues/new/choose)

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@@ -18,12 +18,12 @@ specific language governing permissions and limitations under the License.
# 🧨 Diffusers
🤗 Diffusers provides pretrained vision and audio diffusion models, and serves as a modular toolbox for inference and training.
🤗 Diffusers provides pretrained vision diffusion models, and serves as a modular toolbox for inference and training.
More precisely, 🤗 Diffusers offers:
- State-of-the-art diffusion pipelines that can be run in inference with just a couple of lines of code (see [**Using Diffusers**](./using-diffusers/conditional_image_generation)) or have a look at [**Pipelines**](#pipelines) to get an overview of all supported pipelines and their corresponding papers.
- Various noise schedulers that can be used interchangeably for the preferred speed vs. quality trade-off in inference. For more information see [**Schedulers**](./api/schedulers/overview).
- Various noise schedulers that can be used interchangeably for the preferred speed vs. quality trade-off in inference. For more information see [**Schedulers**](./api/schedulers).
- Multiple types of models, such as UNet, can be used as building blocks in an end-to-end diffusion system. See [**Models**](./api/models) for more details
- Training examples to show how to train the most popular diffusion model tasks. For more information see [**Training**](./training/overview).
@@ -35,7 +35,6 @@ available a colab notebook to directly try them out.
| Pipeline | Paper | Tasks | Colab
|---|---|:---:|:---:|
| [alt_diffusion](./api/pipelines/alt_diffusion) | [**AltDiffusion**](https://arxiv.org/abs/2211.06679) | Image-to-Image Text-Guided Generation |
| [audio_diffusion](./api/pipelines/audio_diffusion) | [**Audio Diffusion**](https://github.com/teticio/audio-diffusion.git) | Unconditional Audio Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/teticio/audio-diffusion/blob/master/notebooks/audio_diffusion_pipeline.ipynb)
| [cycle_diffusion](./api/pipelines/cycle_diffusion) | [**Cycle Diffusion**](https://arxiv.org/abs/2210.05559) | Image-to-Image Text-Guided Generation |
| [dance_diffusion](./api/pipelines/dance_diffusion) | [**Dance Diffusion**](https://github.com/williamberman/diffusers.git) | Unconditional Audio Generation |
| [ddpm](./api/pipelines/ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) | Unconditional Image Generation |
@@ -43,19 +42,14 @@ available a colab notebook to directly try them out.
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation |
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image |
| [latent_diffusion_uncond](./api/pipelines/latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation |
| [paint_by_example](./api/pipelines/paint_by_example) | [**Paint by Example: Exemplar-based Image Editing with Diffusion Models**](https://arxiv.org/abs/2211.13227) | Image-Guided Image Inpainting |
| [pndm](./api/pipelines/pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation |
| [score_sde_ve](./api/pipelines/score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [score_sde_vp](./api/pipelines/score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [stable_diffusion](./api/pipelines/stable_diffusion/text2img) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
| [stable_diffusion](./api/pipelines/stable_diffusion/img2img) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
| [stable_diffusion](./api/pipelines/stable_diffusion/inpaint) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-to-Image Generation |
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting |
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
| [stable_diffusion_safe](./api/pipelines/stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
| [stochastic_karras_ve](./api/pipelines/stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
| [unclip](./api/pipelines/unclip) | [Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) | Text-to-Image Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |

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@@ -120,25 +120,3 @@ git pull
```
Your Python environment will find the `main` version of 🤗 Diffusers on the next run.
## Notice on telemetry logging
Our library gathers telemetry information during `from_pretrained()` requests.
This data includes the version of Diffusers and PyTorch/Flax, the requested model or pipeline class,
and the path to a pretrained checkpoint if it is hosted on the Hub.
This usage data helps us debug issues and prioritize new features.
Telemetry is only sent when loading models and pipelines from the HuggingFace Hub,
and is not collected during local usage.
We understand that not everyone wants to share additional information, and we respect your privacy,
so you can disable telemetry collection by setting the `DISABLE_TELEMETRY` environment variable from your terminal:
On Linux/MacOS:
```bash
export DISABLE_TELEMETRY=YES
```
On Windows:
```bash
set DISABLE_TELEMETRY=YES
```

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@@ -1,193 +0,0 @@
- sections:
- local: index
title: "🧨 Diffusers"
- local: quicktour
title: "훑어보기"
- local: installation
title: "설치"
title: "시작하기"
- sections:
- sections:
- local: in_translation
title: "Loading Pipelines, Models, and Schedulers"
- local: in_translation
title: "Using different Schedulers"
- local: in_translation
title: "Configuring Pipelines, Models, and Schedulers"
- local: in_translation
title: "Loading and Adding Custom Pipelines"
title: "불러오기 & 허브 (번역 예정)"
- sections:
- local: in_translation
title: "Unconditional Image Generation"
- local: in_translation
title: "Text-to-Image Generation"
- local: in_translation
title: "Text-Guided Image-to-Image"
- local: in_translation
title: "Text-Guided Image-Inpainting"
- local: in_translation
title: "Text-Guided Depth-to-Image"
- local: in_translation
title: "Reusing seeds for deterministic generation"
- local: in_translation
title: "Community Pipelines"
- local: in_translation
title: "How to contribute a Pipeline"
title: "추론을 위한 파이프라인 (번역 예정)"
- sections:
- local: in_translation
title: "Reinforcement Learning"
- local: in_translation
title: "Audio"
- local: in_translation
title: "Other Modalities"
title: "Taking Diffusers Beyond Images"
title: "Diffusers 사용법 (번역 예정)"
- sections:
- local: in_translation
title: "Memory and Speed"
- local: in_translation
title: "xFormers"
- local: in_translation
title: "ONNX"
- local: in_translation
title: "OpenVINO"
- local: in_translation
title: "MPS"
- local: in_translation
title: "Habana Gaudi"
title: "최적화/특수 하드웨어 (번역 예정)"
- sections:
- local: in_translation
title: "Overview"
- local: in_translation
title: "Unconditional Image Generation"
- local: in_translation
title: "Textual Inversion"
- local: in_translation
title: "Dreambooth"
- local: in_translation
title: "Text-to-image fine-tuning"
title: "학습 (번역 예정)"
- sections:
- local: in_translation
title: "Stable Diffusion"
- local: in_translation
title: "Philosophy"
- local: in_translation
title: "How to contribute?"
title: "개념 설명 (번역 예정)"
- sections:
- sections:
- local: in_translation
title: "Models"
- local: in_translation
title: "Diffusion Pipeline"
- local: in_translation
title: "Logging"
- local: in_translation
title: "Configuration"
- local: in_translation
title: "Outputs"
title: "Main Classes"
- sections:
- local: in_translation
title: "Overview"
- local: in_translation
title: "AltDiffusion"
- local: in_translation
title: "Cycle Diffusion"
- local: in_translation
title: "DDIM"
- local: in_translation
title: "DDPM"
- local: in_translation
title: "Latent Diffusion"
- local: in_translation
title: "Unconditional Latent Diffusion"
- local: in_translation
title: "PaintByExample"
- local: in_translation
title: "PNDM"
- local: in_translation
title: "Score SDE VE"
- sections:
- local: in_translation
title: "Overview"
- local: in_translation
title: "Text-to-Image"
- local: in_translation
title: "Image-to-Image"
- local: in_translation
title: "Inpaint"
- local: in_translation
title: "Depth-to-Image"
- local: in_translation
title: "Image-Variation"
- local: in_translation
title: "Super-Resolution"
title: "Stable Diffusion"
- local: in_translation
title: "Stable Diffusion 2"
- local: in_translation
title: "Safe Stable Diffusion"
- local: in_translation
title: "Stochastic Karras VE"
- local: in_translation
title: "Dance Diffusion"
- local: in_translation
title: "UnCLIP"
- local: in_translation
title: "Versatile Diffusion"
- local: in_translation
title: "VQ Diffusion"
- local: in_translation
title: "RePaint"
- local: in_translation
title: "Audio Diffusion"
title: "파이프라인 (번역 예정)"
- sections:
- local: in_translation
title: "Overview"
- local: in_translation
title: "DDIM"
- local: in_translation
title: "DDPM"
- local: in_translation
title: "Singlestep DPM-Solver"
- local: in_translation
title: "Multistep DPM-Solver"
- local: in_translation
title: "Heun Scheduler"
- local: in_translation
title: "DPM Discrete Scheduler"
- local: in_translation
title: "DPM Discrete Scheduler with ancestral sampling"
- local: in_translation
title: "Stochastic Kerras VE"
- local: in_translation
title: "Linear Multistep"
- local: in_translation
title: "PNDM"
- local: in_translation
title: "VE-SDE"
- local: in_translation
title: "IPNDM"
- local: in_translation
title: "VP-SDE"
- local: in_translation
title: "Euler scheduler"
- local: in_translation
title: "Euler Ancestral Scheduler"
- local: in_translation
title: "VQDiffusionScheduler"
- local: in_translation
title: "RePaint Scheduler"
title: "스케줄러 (번역 예정)"
- sections:
- local: in_translation
title: "RL Planning"
title: "Experimental Features"
title: "API (번역 예정)"

View File

@@ -1,63 +0,0 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
<p align="center">
<br>
<img src="https://raw.githubusercontent.com/huggingface/diffusers/77aadfee6a891ab9fcfb780f87c693f7a5beeb8e/docs/source/imgs/diffusers_library.jpg" width="400"/>
<br>
</p>
# 🧨 Diffusers
🤗 Diffusers는 사전학습된 비전 및 오디오 확산 모델을 제공하고, 추론 및 학습을 위한 모듈식 도구 상자 역할을 합니다.
보다 정확하게, 🤗 Diffusers는 다음을 제공합니다:
- 단 몇 줄의 코드로 추론을 실행할 수 있는 최신 확산 파이프라인을 제공합니다. ([**Using Diffusers**](./using-diffusers/conditional_image_generation)를 살펴보세요) 지원되는 모든 파이프라인과 해당 논문에 대한 개요를 보려면 [**Pipelines**](#pipelines)을 살펴보세요.
- 추론에서 속도 vs 품질의 절충을 위해 상호교환적으로 사용할 수 있는 다양한 노이즈 스케줄러를 제공합니다. 자세한 내용은 [**Schedulers**](./api/schedulers/overview)를 참고하세요.
- UNet과 같은 여러 유형의 모델을 end-to-end 확산 시스템의 구성 요소로 사용할 수 있습니다. 자세한 내용은 [**Models**](./api/models)을 참고하세요.
- 가장 인기있는 확산 모델 테스크를 학습하는 방법을 보여주는 예제들을 제공합니다. 자세한 내용은 [**Training**](./training/overview)를 참고하세요.
## 🧨 Diffusers 파이프라인
다음 표에는 공시적으로 지원되는 모든 파이프라인, 관련 논문, 직접 사용해 볼 수 있는 Colab 노트북(사용 가능한 경우)이 요약되어 있습니다.
| Pipeline | Paper | Tasks | Colab
|---|---|:---:|:---:|
| [alt_diffusion](./api/pipelines/alt_diffusion) | [**AltDiffusion**](https://arxiv.org/abs/2211.06679) | Image-to-Image Text-Guided Generation |
| [audio_diffusion](./api/pipelines/audio_diffusion) | [**Audio Diffusion**](https://github.com/teticio/audio-diffusion.git) | Unconditional Audio Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/teticio/audio-diffusion/blob/master/notebooks/audio_diffusion_pipeline.ipynb)
| [cycle_diffusion](./api/pipelines/cycle_diffusion) | [**Cycle Diffusion**](https://arxiv.org/abs/2210.05559) | Image-to-Image Text-Guided Generation |
| [dance_diffusion](./api/pipelines/dance_diffusion) | [**Dance Diffusion**](https://github.com/williamberman/diffusers.git) | Unconditional Audio Generation |
| [ddpm](./api/pipelines/ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) | Unconditional Image Generation |
| [ddim](./api/pipelines/ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) | Unconditional Image Generation |
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation |
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image |
| [latent_diffusion_uncond](./api/pipelines/latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation |
| [paint_by_example](./api/pipelines/paint_by_example) | [**Paint by Example: Exemplar-based Image Editing with Diffusion Models**](https://arxiv.org/abs/2211.13227) | Image-Guided Image Inpainting |
| [pndm](./api/pipelines/pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation |
| [score_sde_ve](./api/pipelines/score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [score_sde_vp](./api/pipelines/score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [stable_diffusion](./api/pipelines/stable_diffusion/text2img) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
| [stable_diffusion](./api/pipelines/stable_diffusion/img2img) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
| [stable_diffusion](./api/pipelines/stable_diffusion/inpaint) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-to-Image Generation |
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting |
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
| [stable_diffusion_safe](./api/pipelines/stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
| [stochastic_karras_ve](./api/pipelines/stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
| [unclip](./api/pipelines/unclip) | [Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) | Text-to-Image Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |
| [vq_diffusion](./api/pipelines/vq_diffusion) | [Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation |
**참고**: 파이프라인은 해당 문서에 설명된 대로 확산 시스템을 사용한 방법에 대한 간단한 예입니다.

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@@ -1,142 +0,0 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# 설치
사용하시는 라이브러리에 맞는 🤗 Diffusers를 설치하세요.
🤗 Diffusers는 Python 3.7+, PyTorch 1.7.0+ 및 flax에서 테스트되었습니다. 사용중인 딥러닝 라이브러리에 대한 아래의 설치 안내를 따르세요.
- [PyTorch 설치 안내](https://pytorch.org/get-started/locally/)
- [Flax 설치 안내](https://flax.readthedocs.io/en/latest/)
## pip를 이용한 설치
[가상 환경](https://docs.python.org/3/library/venv.html)에 🤗 Diffusers를 설치해야 합니다.
Python 가상 환경에 익숙하지 않은 경우 [가상환경 pip 설치 가이드](https://packaging.python.org/guides/installing-using-pip-and-virtual-environments/)를 살펴보세요.
가상 환경을 사용하면 서로 다른 프로젝트를 더 쉽게 관리하고, 종속성간의 호환성 문제를 피할 수 있습니다.
프로젝트 디렉토리에 가상 환경을 생성하는 것으로 시작하세요:
```bash
python -m venv .env
```
그리고 가상 환경을 활성화합니다:
```bash
source .env/bin/activate
```
이제 다음의 명령어로 🤗 Diffusers를 설치할 준비가 되었습니다:
**PyTorch의 경우**
```bash
pip install diffusers["torch"]
```
**Flax의 경우**
```bash
pip install diffusers["flax"]
```
## 소스로부터 설치
소스에서 `diffusers`를 설치하기 전에, `torch` 및 `accelerate`이 설치되어 있는지 확인하세요.
`torch` 설치에 대해서는 [torch docs](https://pytorch.org/get-started/locally/#start-locally)를 참고하세요.
다음과 같이 `accelerate`을 설치하세요.
```bash
pip install accelerate
```
다음 명령어를 사용하여 소스에서 🤗 Diffusers를 설치하세요:
```bash
pip install git+https://github.com/huggingface/diffusers
```
이 명령어는 최신 `stable` 버전이 아닌 최첨단 `main` 버전을 설치합니다.
`main` 버전은 최신 개발 정보를 최신 상태로 유지하는 데 유용합니다.
예를 들어 마지막 공식 릴리즈 이후 버그가 수정되었지만, 새 릴리즈가 아직 출시되지 않은 경우입니다.
그러나 이는 `main` 버전이 항상 안정적이지 않을 수 있음을 의미합니다.
우리는 `main` 버전이 지속적으로 작동하도록 노력하고 있으며, 대부분의 문제는 보통 몇 시간 또는 하루 안에 해결됩니다.
문제가 발생하면 더 빨리 해결할 수 있도록 [Issue](https://github.com/huggingface/transformers/issues)를 열어주세요!
## 편집가능한 설치
다음을 수행하려면 편집가능한 설치가 필요합니다:
* 소스 코드의 `main` 버전을 사용
* 🤗 Diffusers에 기여 (코드의 변경 사항을 테스트하기 위해 필요)
저장소를 복제하고 다음 명령어를 사용하여 🤗 Diffusers를 설치합니다:
```bash
git clone https://github.com/huggingface/diffusers.git
cd diffusers
```
**PyTorch의 경우**
```
pip install -e ".[torch]"
```
**Flax의 경우**
```
pip install -e ".[flax]"
```
이러한 명령어들은 저장소를 복제한 폴더와 Python 라이브러리 경로를 연결합니다.
Python은 이제 일반 라이브러리 경로에 더하여 복제한 폴더 내부를 살펴봅니다.
예를들어 Python 패키지가 `~/anaconda3/envs/main/lib/python3.7/site-packages/`에 설치되어 있는 경우 Python은 복제한 폴더인 `~/diffusers/`도 검색합니다.
<Tip warning={true}>
라이브러리를 계속 사용하려면 `diffusers` 폴더를 유지해야 합니다.
</Tip>
이제 다음 명령어를 사용하여 최신 버전의 🤗 Diffusers로 쉽게 업데이트할 수 있습니다:
```bash
cd ~/diffusers/
git pull
```
이렇게 하면, 다음에 실행할 때 Python 환경이 🤗 Diffusers의 `main` 버전을 찾게 됩니다.
## 텔레메트리 로깅에 대한 알림
우리 라이브러리는 `from_pretrained()` 요청 중에 텔레메트리 정보를 원격으로 수집합니다.
이 데이터에는 Diffusers 및 PyTorch/Flax의 버전, 요청된 모델 또는 파이프라인 클래스, 그리고 허브에서 호스팅되는 경우 사전학습된 체크포인트에 대한 경로를 포함합니다.
이 사용 데이터는 문제를 디버깅하고 새로운 기능의 우선순위를 지정하는데 도움이 됩니다.
텔레메트리는 HuggingFace 허브에서 모델과 파이프라인을 불러올 때만 전송되며, 로컬 사용 중에는 수집되지 않습니다.
우리는 추가 정보를 공유하지 않기를 원하는 사람이 있다는 것을 이해하고 개인 정보를 존중하므로, 터미널에서 `DISABLE_TELEMETRY` 환경 변수를 설정하여 텔레메트리 수집을 비활성화할 수 있습니다.
Linux/MacOS에서:
```bash
export DISABLE_TELEMETRY=YES
```
Windows에서:
```bash
set DISABLE_TELEMETRY=YES
```

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@@ -1,123 +0,0 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# 훑어보기
🧨 Diffusers로 빠르게 시작하고 실행하세요!
이 훑어보기는 여러분이 개발자, 일반사용자 상관없이 시작하는 데 도움을 주며, 추론을 위해 [`DiffusionPipeline`] 사용하는 방법을 보여줍니다.
시작하기에 앞서서, 필요한 모든 라이브러리가 설치되어 있는지 확인하세요:
```bash
pip install --upgrade diffusers accelerate transformers
```
- [`accelerate`](https://huggingface.co/docs/accelerate/index)은 추론 및 학습을 위한 모델 불러오기 속도를 높입니다.
- [`transformers`](https://huggingface.co/docs/transformers/index)는 [Stable Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/overview)과 같이 가장 널리 사용되는 확산 모델을 실행하기 위해 필요합니다.
## DiffusionPipeline
[`DiffusionPipeline`]은 추론을 위해 사전학습된 확산 시스템을 사용하는 가장 쉬운 방법입니다. 다양한 양식의 많은 작업에 [`DiffusionPipeline`]을 바로 사용할 수 있습니다. 지원되는 작업은 아래의 표를 참고하세요:
| **Task** | **Description** | **Pipeline**
|------------------------------|--------------------------------------------------------------------------------------------------------------|-----------------|
| Unconditional Image Generation | 가우시안 노이즈에서 이미지 생성 | [unconditional_image_generation](./using-diffusers/unconditional_image_generation`) |
| Text-Guided Image Generation | 텍스트 프롬프트로 이미지 생성 | [conditional_image_generation](./using-diffusers/conditional_image_generation) |
| Text-Guided Image-to-Image Translation | 텍스트 프롬프트에 따라 이미지 조정 | [img2img](./using-diffusers/img2img) |
| Text-Guided Image-Inpainting | 마스크 및 텍스트 프롬프트가 주어진 이미지의 마스킹된 부분을 채우기 | [inpaint](./using-diffusers/inpaint) |
| Text-Guided Depth-to-Image Translation | 깊이 추정을 통해 구조를 유지하면서 텍스트 프롬프트에 따라 이미지의 일부를 조정 | [depth2image](./using-diffusers/depth2image) |
확산 파이프라인이 다양한 작업에 대해 어떻게 작동하는지는 [**Using Diffusers**](./using-diffusers/overview)를 참고하세요.
예를들어, [`DiffusionPipeline`] 인스턴스를 생성하여 시작하고, 다운로드하려는 파이프라인 체크포인트를 지정합니다.
모든 [Diffusers' checkpoint](https://huggingface.co/models?library=diffusers&sort=downloads)에 대해 [`DiffusionPipeline`]을 사용할 수 있습니다.
하지만, 이 가이드에서는 [Stable Diffusion](https://huggingface.co/CompVis/stable-diffusion)을 사용하여 text-to-image를 하는데 [`DiffusionPipeline`]을 사용합니다.
[Stable Diffusion](https://huggingface.co/CompVis/stable-diffusion) 기반 모델을 실행하기 전에 [license](https://huggingface.co/spaces/CompVis/stable-diffusion-license)를 주의 깊게 읽으세요.
이는 모델의 향상된 이미지 생성 기능과 이것으로 생성될 수 있는 유해한 콘텐츠 때문입니다. 선택한 Stable Diffusion 모델(*예*: [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5))로 이동하여 라이센스를 읽으세요.
다음과 같이 모델을 로드할 수 있습니다:
```python
>>> from diffusers import DiffusionPipeline
>>> pipeline = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
```
[`DiffusionPipeline`]은 모든 모델링, 토큰화 및 스케줄링 구성요소를 다운로드하고 캐시합니다.
모델은 약 14억개의 매개변수로 구성되어 있으므로 GPU에서 실행하는 것이 좋습니다.
PyTorch에서와 마찬가지로 생성기 객체를 GPU로 옮길 수 있습니다.
```python
>>> pipeline.to("cuda")
```
이제 `pipeline`을 사용할 수 있습니다:
```python
>>> image = pipeline("An image of a squirrel in Picasso style").images[0]
```
출력은 기본적으로 [PIL Image object](https://pillow.readthedocs.io/en/stable/reference/Image.html?highlight=image#the-image-class)로 래핑됩니다.
다음과 같이 함수를 호출하여 이미지를 저장할 수 있습니다:
```python
>>> image.save("image_of_squirrel_painting.png")
```
**참고**: 다음을 통해 가중치를 다운로드하여 로컬에서 파이프라인을 사용할 수도 있습니다:
```
git lfs install
git clone https://huggingface.co/runwayml/stable-diffusion-v1-5
```
그리고 저장된 가중치를 파이프라인에 불러옵니다.
```python
>>> pipeline = DiffusionPipeline.from_pretrained("./stable-diffusion-v1-5")
```
파이프라인 실행은 동일한 모델 아키텍처이므로 위의 코드와 동일합니다.
```python
>>> generator.to("cuda")
>>> image = generator("An image of a squirrel in Picasso style").images[0]
>>> image.save("image_of_squirrel_painting.png")
```
확산 시스템은 각각 장점이 있는 여러 다른 [schedulers](./api/schedulers/overview)와 함께 사용할 수 있습니다. 기본적으로 Stable Diffusion은 `PNDMScheduler`로 실행되지만 다른 스케줄러를 사용하는 방법은 매우 간단합니다. *예* [`EulerDiscreteScheduler`] 스케줄러를 사용하려는 경우, 다음과 같이 사용할 수 있습니다:
```python
>>> from diffusers import EulerDiscreteScheduler
>>> pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
>>> # change scheduler to Euler
>>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
```
스케줄러 변경 방법에 대한 자세한 내용은 [Using Schedulers](./using-diffusers/schedulers) 가이드를 참고하세요.
[Stability AI's](https://stability.ai/)의 Stable Diffusion 모델은 인상적인 이미지 생성 모델이며 텍스트에서 이미지를 생성하는 것보다 훨씬 더 많은 작업을 수행할 수 있습니다. 우리는 Stable Diffusion만을 위한 전체 문서 페이지를 제공합니다 [link](./conceptual/stable_diffusion).
만약 더 적은 메모리, 더 높은 추론 속도, Mac과 같은 특정 하드웨어 또는 ONNX 런타임에서 실행되도록 Stable Diffusion을 최적화하는 방법을 알고 싶다면 최적화 페이지를 살펴보세요:
- [Optimized PyTorch on GPU](./optimization/fp16)
- [Mac OS with PyTorch](./optimization/mps)
- [ONNX](./optimization/onnx)
- [OpenVINO](./optimization/open_vino)
확산 모델을 미세조정하거나 학습시키려면, [**training section**](./training/overview)을 살펴보세요.
마지막으로, 생성된 이미지를 공개적으로 배포할 때 신중을 기해 주세요 🤗.

View File

@@ -12,14 +12,13 @@ specific language governing permissions and limitations under the License.
# Memory and speed
We present some techniques and ideas to optimize 🤗 Diffusers _inference_ for memory or speed. As a general rule, we recommend the use of [xFormers](https://github.com/facebookresearch/xformers) for memory efficient attention, please see the recommended [installation instructions](xformers).
We'll discuss how the following settings impact performance and memory.
We present some techniques and ideas to optimize 🤗 Diffusers _inference_ for memory or speed.
| | Latency | Speedup |
| ---------------- | ------- | ------- |
| original | 9.50s | x1 |
| cuDNN auto-tuner | 9.37s | x1.01 |
| autocast (fp16) | 5.47s | x1.74 |
| fp16 | 3.61s | x2.63 |
| channels last | 3.30s | x2.88 |
| traced UNet | 3.21s | x2.96 |
@@ -53,14 +52,32 @@ import torch
torch.backends.cuda.matmul.allow_tf32 = True
```
## Automatic mixed precision (AMP)
If you use a CUDA GPU, you can take advantage of `torch.autocast` to perform inference roughly twice as fast at the cost of slightly lower precision. All you need to do is put your inference call inside an `autocast` context manager. The following example shows how to do it using Stable Diffusion text-to-image generation as an example:
```Python
from torch import autocast
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
with autocast("cuda"):
image = pipe(prompt).images[0]
```
Despite the precision loss, in our experience the final image results look the same as the `float32` versions. Feel free to experiment and report back!
## Half precision weights
To save more GPU memory and get more speed, you can load and run the model weights directly in half precision. This involves loading the float16 version of the weights, which was saved to a branch named `fp16`, and telling PyTorch to use the `float16` type when loading them:
To save more GPU memory and get even more speed, you can load and run the model weights directly in half precision. This involves loading the float16 version of the weights, which was saved to a branch named `fp16`, and telling PyTorch to use the `float16` type when loading them:
```Python
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
@@ -69,11 +86,6 @@ prompt = "a photo of an astronaut riding a horse on mars"
image = pipe(prompt).images[0]
```
<Tip warning={true}>
It is strongly discouraged to make use of [`torch.autocast`](https://pytorch.org/docs/stable/amp.html#torch.autocast) in any of the pipelines as it can lead to black images and is always slower than using pure
float16 precision.
</Tip>
## Sliced attention for additional memory savings
For even additional memory savings, you can use a sliced version of attention that performs the computation in steps instead of all at once.
@@ -93,7 +105,7 @@ from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
@@ -105,37 +117,9 @@ image = pipe(prompt).images[0]
There's a small performance penalty of about 10% slower inference times, but this method allows you to use Stable Diffusion in as little as 3.2 GB of VRAM!
## Sliced VAE decode for larger batches
To decode large batches of images with limited VRAM, or to enable batches with 32 images or more, you can use sliced VAE decode that decodes the batch latents one image at a time.
You likely want to couple this with [`~StableDiffusionPipeline.enable_attention_slicing`] or [`~StableDiffusionPipeline.enable_xformers_memory_efficient_attention`] to further minimize memory use.
To perform the VAE decode one image at a time, invoke [`~StableDiffusionPipeline.enable_vae_slicing`] in your pipeline before inference. For example:
```Python
import torch
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
pipe.enable_vae_slicing()
images = pipe([prompt] * 32).images
```
You may see a small performance boost in VAE decode on multi-image batches. There should be no performance impact on single-image batches.
## Offloading to CPU with accelerate for memory savings
For additional memory savings, you can offload the weights to CPU and only load them to GPU when performing the forward pass.
For additional memory savings, you can offload the weights to CPU and load them to GPU when performing the forward pass.
To perform CPU offloading, all you have to do is invoke [`~StableDiffusionPipeline.enable_sequential_cpu_offload`]:
@@ -145,18 +129,19 @@ from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
pipe.enable_sequential_cpu_offload()
image = pipe(prompt).images[0]
```
And you can get the memory consumption to < 3GB.
And you can get the memory consumption to < 2GB.
If is also possible to chain it with attention slicing for minimal memory consumption (< 2GB).
If is also possible to chain it with attention slicing for minimal memory consumption, running it in as little as < 800mb of GPU vRAM:
```Python
import torch
@@ -164,9 +149,10 @@ from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
pipe.enable_sequential_cpu_offload()
@@ -175,8 +161,6 @@ pipe.enable_attention_slicing(1)
image = pipe(prompt).images[0]
```
**Note**: When using `enable_sequential_cpu_offload()`, it is important to **not** move the pipeline to CUDA beforehand or else the gain in memory consumption will only be minimal. See [this issue](https://github.com/huggingface/diffusers/issues/1934) for more information.
## Using Channels Last memory format
Channels last memory format is an alternative way of ordering NCHW tensors in memory preserving dimensions ordering. Channels last tensors ordered in such a way that channels become the densest dimension (aka storing images pixel-per-pixel). Since not all operators currently support channels last format it may result in a worst performance, so it's better to try it and see if it works for your model.
@@ -220,6 +204,7 @@ def generate_inputs():
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
).to("cuda")
unet = pipe.unet
@@ -283,6 +268,7 @@ class UNet2DConditionOutput:
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
).to("cuda")
@@ -308,9 +294,7 @@ with torch.inference_mode():
## Memory Efficient Attention
Recent work on optimizing the bandwitdh in the attention block has generated huge speed ups and gains in GPU memory usage. The most recent being Flash Attention from @tridao: [code](https://github.com/HazyResearch/flash-attention), [paper](https://arxiv.org/pdf/2205.14135.pdf).
Recent work on optimizing the bandwitdh in the attention block have generated huge speed ups and gains in GPU memory usage. The most recent being Flash Attention (from @tridao, [code](https://github.com/HazyResearch/flash-attention), [paper](https://arxiv.org/pdf/2205.14135.pdf)) .
Here are the speedups we obtain on a few Nvidia GPUs when running the inference at 512x512 with a batch size of 1 (one prompt):
| GPU | Base Attention FP16 | Memory Efficient Attention FP16 |
@@ -326,13 +310,14 @@ Here are the speedups we obtain on a few Nvidia GPUs when running the inference
To leverage it just make sure you have:
- PyTorch > 1.12
- Cuda available
- [Installed the xformers library](xformers).
- Installed the [xformers](https://github.com/facebookresearch/xformers) library
```python
from diffusers import StableDiffusionPipeline
import torch
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
).to("cuda")
@@ -343,4 +328,4 @@ with torch.inference_mode():
# optional: You can disable it via
# pipe.disable_xformers_memory_efficient_attention()
```
```

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