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145 Commits

Author SHA1 Message Date
anton-
f7bb9cae11 Release: v0.11.1 2022-12-20 17:24:39 +01:00
Suraj Patil
40b0519a8a Fix num images per prompt unclip (#1787)
* use repeat_interleave

* fix repeat

* Trigger Build

* don't install accelerate from main

* install released accelrate for mps test

* Remove additional accelerate installation from main.

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-20 17:23:11 +01:00
anton-
a5edb981a7 [Patch] Return import for the unclip pipeline loader 2022-12-19 17:56:42 +01:00
anton-
54796b7e43 Release: v0.11.0 2022-12-19 17:43:22 +01:00
Anton Lozhkov
4cb887e0a7 Transformers version req for UnCLIP (#1766)
* Transformers version req for UnCLIP

* add to the list
2022-12-19 17:11:17 +01:00
Anish Shah
9f657f106d [Examples] Update train_unconditional.py to include logging argument for Wandb (#1719)
Update train_unconditional.py

Add logger flag to choose between tensorboard and wandb
2022-12-19 16:57:03 +01:00
Patrick von Platen
ce1c27adc8 [Revision] Don't recommend using revision (#1764) 2022-12-19 16:25:41 +01:00
Patrick von Platen
b267d28566 [Versatile] fix attention mask (#1763) 2022-12-19 15:58:39 +01:00
Anton Lozhkov
c7b4acfb37 Add CPU offloading to UnCLIP (#1761)
* Add CPU offloading to UnCLIP

* use fp32 for testing the offload
2022-12-19 14:44:08 +01:00
Suraj Patil
be38b2d711 [UnCLIPPipeline] fix num_images_per_prompt (#1762)
duplicate maks for num_images_per_prompt
2022-12-19 14:32:46 +01:00
Anton Lozhkov
32a5d70c42 Support attn2==None for xformers (#1759) 2022-12-19 12:43:30 +01:00
Patrick von Platen
429e5449c1 Add attention mask to uclip (#1756)
* Remove bogus file

* [Unclip] Add efficient attention

* [Unclip] Add efficient attention
2022-12-19 12:10:46 +01:00
Anton Lozhkov
dc7cd893fd Add resnet_time_scale_shift to VD layers (#1757) 2022-12-19 12:01:46 +01:00
Mikołaj Siedlarek
8890758823 Correct help text for scheduler_type flag in scripts. (#1749) 2022-12-19 11:27:23 +01:00
Will Berman
b25843e799 unCLIP docs (#1754)
* [unCLIP docs] markdown

* [unCLIP docs] UnCLIPPipeline
2022-12-19 10:27:32 +01:00
Will Berman
830a9d1f01 [fix] pipeline_unclip generator (#1751)
* [fix] pipeline_unclip generator

pass generator to all schedulers

* fix fast tests test data
2022-12-19 10:27:18 +01:00
Will Berman
2dcf64b72a kakaobrain unCLIP (#1428)
* [wip] attention block updates

* [wip] unCLIP unet decoder and super res

* [wip] unCLIP prior transformer

* [wip] scheduler changes

* [wip] text proj utility class

* [wip] UnCLIPPipeline

* [wip] kakaobrain unCLIP convert script

* [unCLIP pipeline] fixes re: @patrickvonplaten

remove callbacks

move denoising loops into call function

* UNCLIPScheduler re: @patrickvonplaten

Revert changes to DDPMScheduler. Make UNCLIPScheduler, a modified
DDPM scheduler with changes to support karlo

* mask -> attention_mask re: @patrickvonplaten

* [DDPMScheduler] remove leftover change

* [docs] PriorTransformer

* [docs] UNet2DConditionModel and UNet2DModel

* [nit] UNCLIPScheduler -> UnCLIPScheduler

matches existing unclip naming better

* [docs] SchedulingUnCLIP

* [docs] UnCLIPTextProjModel

* refactor

* finish licenses

* rename all to attention_mask and prep in models

* more renaming

* don't expose unused configs

* final renaming fixes

* remove x attn mask when not necessary

* configure kakao script to use new class embedding config

* fix copies

* [tests] UnCLIPScheduler

* finish x attn

* finish

* remove more

* rename condition blocks

* clean more

* Apply suggestions from code review

* up

* fix

* [tests] UnCLIPPipelineFastTests

* remove unused imports

* [tests] UnCLIPPipelineIntegrationTests

* correct

* make style

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-18 15:15:30 -08:00
Patrick von Platen
402b9560b2 Remove license accept ticks 2022-12-19 00:10:17 +01:00
Anton Lozhkov
c2a38ef9df Fix/update the LDM pipeline and tests (#1743)
* Fix/update LDM tests

* batched generators
2022-12-18 11:49:53 +01:00
Anton Lozhkov
08cc36ddff Fix MPS fast test warnings (#1744)
* unset level
2022-12-17 22:57:30 +01:00
Peter
723e8f6bb4 Fix ONNX img2img preprocessing (#1736)
Co-authored-by: Peter <peterto@users.noreply.github.com>
2022-12-17 13:12:10 +01:00
Patrick von Platen
c53a850604 [Batched Generators] This PR adds generators that are useful to make batched generation fully reproducible (#1718)
* [Batched Generators] all batched generators

* up

* up

* up

* up

* up

* up

* up

* up

* up

* up

* up

* up

* up

* up

* up

* up

* hey

* up again

* fix tests

* Apply suggestions from code review

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

* correct tests

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-17 11:13:16 +01:00
Anton Lozhkov
086c7f9ea8 Nightly integration tests (#1664)
* [WIP] Nightly integration tests

* initial SD tests

* update SD slow tests

* style

* repaint

* ImageVariations

* style

* finish imgvar

* img2img tests

* debug

* inpaint 1.5

* inpaint legacy

* torch isn't happy about deterministic ops

* allclose -> max diff for shorter logs

* add SD2

* debug

* Update tests/pipelines/stable_diffusion_2/test_stable_diffusion.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* Update tests/pipelines/stable_diffusion/test_stable_diffusion.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* fix refs

* Update src/diffusers/utils/testing_utils.py

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

* fix refs

* remove debug

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-16 18:51:11 +01:00
Pedro Cuenca
acd317810b Docs: recommend xformers (#1724)
* Fix links to flash attention.

* Add xformers installation instructions.

* Make link to xformers install more prominent.

* Link to xformers install from training docs.
2022-12-16 15:49:01 +01:00
Patrick von Platen
c6d0dff4a3 Fix ldm tests on master by not running the CPU tests on GPU (#1729) 2022-12-16 15:28:40 +01:00
Anton Lozhkov
a40095dd22 Fix ONNX img2img preprocessing and add fast tests coverage (#1727)
* Fix ONNX img2img preprocessing and add fast tests coverage

* revert

* disable progressbars
2022-12-16 15:24:16 +01:00
Partho
727434c206 Accept latents as optional input in Latent Diffusion pipeline (#1723)
* Latent Diffusion pipeline accept latents

* make style

* check for mps

randn does not work reproducibly on mps
2022-12-16 12:13:41 +01:00
YiYi Xu
21e61eb3a9 Added a README page for docs and a "schedulers" page (#1710)
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-15 13:04:40 -10:00
Haihao Shen
c891330f79 Add examples with Intel optimizations (#1579)
* Add examples with Intel optimizations (BF16 fine-tuning and inference)

* Remove unused package

* Add README for intel_opts and refine the description for research projects

* Add notes of intel opts for diffusers
2022-12-15 21:16:27 +01:00
jiqing-feng
c5f04d4e34 apply amp bf16 on textual inversion (#1465)
* add conf.yaml

* enable bf16

enable amp bf16 for unet forward

fix style

fix readme

remove useless file

* change amp to full bf16

* align

* make stype

* fix format
2022-12-15 21:15:23 +01:00
CyberMeow
61dec53356 Improve pipeline_stable_diffusion_inpaint_legacy.py (#1585)
* update inpaint_legacy to allow the use of predicted noise to construct intermediate diffused images

* Update src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_inpaint_legacy.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-15 20:59:31 +01:00
Pedro Cuenca
badddee0ef Add state checkpointing to other training scripts (#1687)
* Add state checkpointing to other training scripts

* Fix first_epoch

* Apply suggestions from code review

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* Update Dreambooth checkpoint help message.

* Dreambooth docs: checkpoints, inference from a checkpoint.

* make style

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-15 19:49:40 +01:00
Anton Lozhkov
13994b2d3f RePaint fast tests and API conforming (#1701)
* add fast tests

* better tests and fp16

* batch fixes

* Reuse preprocessing

* quickfix
2022-12-15 18:35:31 +01:00
anton-
ea90bf2ba1 skip mps 2022-12-15 18:01:39 +01:00
Chino
8cecc66a74 Fix the bug that torch version less than 1.12 throws TypeError (#1671) 2022-12-14 21:29:39 +01:00
Anton Lozhkov
35b66c8e32 [Readme] Clarify package owners (#1707)
Specify that we don't actively monitor the conda scripts
2022-12-14 20:49:36 +01:00
Patrick von Platen
013edb641a Update main docs (#1706)
* Remove bogus file

* [Docs] Remove mentioning of gated access since no longer exsits

* add docs to index

* Apply suggestions from code review

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-14 20:33:54 +01:00
Anton Lozhkov
86ac3ea1d7 Delete _ 2022-12-14 13:52:29 +01:00
Anton Lozhkov
ef3fcbb688 Remove all local telemetry (#1702) 2022-12-14 12:56:35 +01:00
Prathik Rao
7c823c2ed7 manually update train_unconditional_ort (#1694)
* manually update train_unconditional_ort

* formatting

Co-authored-by: Prathik Rao <prathikrao@microsoft.com>
2022-12-14 11:35:41 +01:00
Pedro Cuenca
784beee969 Dreambooth: use warnings instead of logger in parse_args() (#1688)
Use warnings instead of logger in parse_args()

logger requires an `Accelerator`.
2022-12-13 22:01:48 +01:00
Patrick von Platen
8b7cb962a5 make style 2022-12-13 17:01:18 +00:00
Patrick von Platen
e1bb8f6188 [Community pipeline] Add github mechanism (#1680)
* [Community pipeline] Add github mechanism

* better

* Apply suggestions from code review

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

* adapt

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-13 18:01:00 +01:00
Patrick von Platen
e62dd5cfa8 Change one-step dummy pipeline for testing (#1690)
Change the one-step dummy pipeline for testing
2022-12-13 16:55:49 +01:00
w4ffl35
07f95503e5 Disable telemetry when DISABLE_TELEMETRY is set (#1686)
fixed #1685 - disables telemetry when DISABLE_TELEMETRY and HF_HUB_OFFLINE is set
2022-12-13 16:28:07 +01:00
Pedro Cuenca
e01d6cf295 Dreambooth: save / restore training state (#1668)
* Dreambooth: save / restore training state.

* make style

* Rename vars for clarity.

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* Remove unused import

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-13 15:16:44 +01:00
Patrick von Platen
244e16a7ab [Version] Bump to 0.11.0.dev0 (#1682)
upgrade version
2022-12-13 13:51:36 +01:00
Patrick von Platen
b345c74d4d Make sure all pipelines can run with batched input (#1669)
* [SD] Make sure batched input works correctly

* uP

* uP

* up

* up

* uP

* up

* fix mask stuff

* up

* uP

* more up

* up

* uP

* up

* finish

* Apply suggestions from code review

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-13 12:50:15 +01:00
apolinario
b417042291 Fix wrong type checking in convert_diffusers_to_original_stable_diffusion.py (#1681)
* Fix type checking remainders

* Remove IS_V20_MODEL flag always being True

Co-authored-by: apolinario <joaopaulo.passos+multimodal@gmail.com>
2022-12-13 12:44:20 +01:00
Suvaditya Mukherjee
40c16ed2f0 Added Community pipeline for comparing Stable Diffusion v1.1-4 checkpoints (#1584)
* Added Community pipeline for comparing Stable Diffusion v1.1-4

Signed-off-by: Suvaditya Mukherjee <suvadityamuk@gmail.com>

* Made changes to provide support for current iteration of from_pretrained and added example

Signed-off-by: Suvaditya Mukherjee <suvadityamuk@gmail.com>

* updated a small spelling error

Signed-off-by: Suvaditya Mukherjee <suvadityamuk@gmail.com>

* added pipeline entry to table

Signed-off-by: Suvaditya Mukherjee <suvadityamuk@gmail.com>

Signed-off-by: Suvaditya Mukherjee <suvadityamuk@gmail.com>
2022-12-13 11:31:30 +01:00
Patrick von Platen
69de9b2eaa [Textual Inversion] Do not update other embeddings (#1665) 2022-12-12 17:44:39 +01:00
Patrick von Platen
3ce6380d3a [SD] Make sure scheduler is correct when converting (#1667) 2022-12-12 16:57:48 +01:00
Cyberes
d2dc4de303 Handle missing global_step key in scripts/convert_original_stable_diffusion_to_diffusers.py (#1612)
handle missing global_step key and don't download config if it already exists
2022-12-12 16:10:52 +01:00
Kangfu Mei
ded3299d68 fix bug if we don't do_classifier_free_guidance (#1601)
* fix bug if we don't do_classifier_free_guidance

* update for copied diffusers.pipelines..alt_diffusion..pipeline_alt_diffusion.AltDiffusionPipeline
2022-12-12 15:02:13 +01:00
Patrick von Platen
8bf5e59931 Deprecate init image correctly (#1649)
Deprecate init image correctl
2022-12-12 15:00:20 +01:00
Prathik Rao
4645e28355 tensor format ort bug fix (#1557)
bug fix

Co-authored-by: Prathik Rao <prathikrao@microsoft.com>
Co-authored-by: anton- <anton@huggingface.co>
2022-12-12 13:56:02 +01:00
Lukas Struppek
589330595d VersatileDiffusion: fix input processing (#1568)
* fix versatile diffusion input

* merge main

* `make fix-copies`

Co-authored-by: anton- <anton@huggingface.co>
2022-12-12 13:45:27 +01:00
lawfordp2017
31444f5790 Add text encoder conversion (#1559)
* Initial code for attempt at improving SD <--> diffusers conversions for v2.0

* Updates to support round-trip between orig. SD 2.0 and diffusers models

* Corrected formatting to Black standard

* Correcting import formatting

* Fixed imports (properly this time)

* add some corrections

* remove inference files

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-12 10:07:42 +01:00
M-A
c3b2f97534 Remove unnecessary kwargs in depth2img (#1648)
Since no deprecate() call was done, the typos were silently ignored.
2022-12-11 17:00:50 +01:00
Patrick von Platen
fc94c60c83 Remove unnecessary offset in img2img (#1653)
remove unnecessary offset in img2img
2022-12-10 19:26:25 +01:00
Pedro Cuenca
ea64a7860a Allow k pipeline to generate > 1 images (#1645)
Allow k pipeline to generate > 1 images.
2022-12-10 17:54:02 +01:00
Tim Hinderliter
2868d99181 dreambooth: fix #1566: maintain fp32 wrapper when saving a checkpoint to avoid crash when running fp16 (#1618)
* dreambooth: fix #1566: maintain fp32 wrapper when saving a checkpoint to avoid crash when running fp16

* dreambooth: guard against passing keep_fp32_wrapper arg to older versions of accelerate. part of fix for #1566

* Apply suggestions from code review

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

* Update examples/dreambooth/train_dreambooth.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-10 15:45:45 +01:00
Pedro Cuenca
0c18d02cc9 Remove spurious arg in training scripts (#1644)
Remove spurious arg in training scripts.
2022-12-10 13:57:20 +01:00
Patrick von Platen
6b68afd8e4 do not automatically enable xformers (#1640)
* do not automatically enable xformers

* uP
2022-12-09 18:28:36 +01:00
anton-
63c4944998 Patch release: v0.10.2 2022-12-09 17:59:32 +01:00
Anton Lozhkov
3ebe40fc5f Adapt to forced transformers version in some dependent libraries (#1638)
* Adapt to forced transformers version in some dependent libraries

* style

* Update __init__.py

* update requires_backends
2022-12-09 17:42:53 +01:00
Anton Lozhkov
089252542c V0.10.1 patch (#1637)
* Re-add xformers enable to UNet2DCondition (#1627)

* finish

* fix

* Update tests/models/test_models_unet_2d.py

* style

Co-authored-by: Anton Lozhkov <anton@huggingface.co>

* Release: v0.10.1

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-09 17:04:25 +01:00
Patrick von Platen
cd91fc06fe Re-add xformers enable to UNet2DCondition (#1627)
* finish

* fix

* Update tests/models/test_models_unet_2d.py

* style

Co-authored-by: Anton Lozhkov <anton@huggingface.co>
2022-12-09 14:05:38 +01:00
Pedro Cuenca
ff65c2d72b Don't assume 512x512 in k-diffusion pipeline (#1625)
Don't assume 512x512 in k-diffusion pipeline.
2022-12-09 11:50:29 +01:00
Haofan Wang
f1b726e46e Update requirements.txt (#1623)
* Update requirements.txt

* Update requirements_flax.txt

* Update requirements.txt

* Update requirements_flax.txt

* Update requirements.txt

* Update requirements_flax.txt
2022-12-09 08:35:27 +01:00
SkyTNT
f242eba4fd Fix lpw stable diffusion pipeline compatibility (#1622) 2022-12-09 08:30:26 +01:00
anton-
3faf204c49 Release: v0.10.0 2022-12-08 19:24:10 +01:00
Suraj Patil
5383188c7e StableDiffusionDepth2ImgPipeline (#1531)
* begin depth pipeline

* add depth estimation model

* fix prepare_depth_mask

* add a comment about autocast

* copied from, quality, cleanup

* begin tests

* handle tensors

* norm image tensor

* fix batch size

* fix tests

* fix enable_sequential_cpu_offload

* fix save load

* fix test_save_load_float16

* fix test_save_load_optional_components

* fix test_float16_inference

* fix test_cpu_offload_forward_pass

* fix test_dict_tuple_outputs_equivalent

* up

* fix fast tests

* fix test_stable_diffusion_img2img_multiple_init_images

* fix few more fast tests

* don't use device map for DPT

* fix test_stable_diffusion_pipeline_with_sequential_cpu_offloading

* accept external depth maps

* prepare_depth_mask -> prepare_depth_map

* fix file name

* fix file name

* quality

* check transformers version

* fix test names

* use skipif

* fix import

* add docs

* skip tests on mps

* correct version

* uP

* Update docs/source/api/pipelines/stable_diffusion_2.mdx

* fix fix-copies

* fix fix-copies

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
Co-authored-by: anton- <anton@huggingface.co>
2022-12-08 18:25:12 +01:00
Anton Lozhkov
dbe0719246 Fix PyCharm/VSCode static type checking for dummy objects (#1596)
* Fix PyCharm/VSCode static type checking for dummy objects

* Re-add dummies

* Fix AudioDiffusion imports

* fix import

* fix import

* Update utils/check_dummies.py

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

* Update src/diffusers/utils/import_utils.py

* Update src/diffusers/__init__.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* Update src/diffusers/pipelines/stable_diffusion/__init__.py

* fix double import

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-08 14:02:11 +01:00
Anton Lozhkov
03566d8689 Delete hi 2022-12-08 13:07:25 +01:00
Suraj Patil
a934e5bc6c [Versatile Diffusion] add upcast_attention (#1605)
add upcast_attention arg
2022-12-08 13:03:32 +01:00
Patrick von Platen
a643c6300e [K Diffusion] Add k diffusion sampler natively (#1603)
* uP

* uP
2022-12-08 12:48:37 +01:00
Ben Sherman
326de41915 Trivial fix for undefined symbol in train_dreambooth.py (#1598)
easy fix for undefined name in train_dreambooth.py

import_model_class_from_model_name_or_path loads a pretrained model
and refers to args.revision in a context where args is undefined. I modified
the function to take revision as an argument and modified the invocation
of the function to pass in the revision from args. Seems like this was caused
by a cut and paste.
2022-12-07 21:39:48 +01:00
Anton Lozhkov
eb1abee693 [ONNX] Fix flaky tests (#1593)
* [ONNX] Fix flaky tests

* revert
2022-12-07 19:53:13 +01:00
Pedro Cuenca
5e0369219f Make cross-attention check more robust (#1560)
* Make cross-attention check more robust.

* Fix copies.
2022-12-07 18:33:29 +01:00
Nathan Lambert
bea7eb4314 Update RL docs for better sharing / adding models (#1563)
* init docs update

* style

* fix bad colab formatting, add pipeline comment

* update todo
2022-12-07 09:08:12 -08:00
Randolph-zeng
ca68ab3eef Update scheduling_repaint.py (#1582)
* Update scheduling_repaint.py

* update the expected image

Co-authored-by: anton- <anton@huggingface.co>
2022-12-07 17:41:07 +01:00
Suraj Patil
ced7c9601a fix upcast in slice attention (#1591)
* fix upcast in slice attention

* fix dtype

* add test

* fix test
2022-12-07 15:14:34 +01:00
Cheng Lu
8e74efad01 Add Singlestep DPM-Solver (singlestep high-order schedulers) (#1442)
* add singlestep dpmsolver

* fix a style typo

* fix a style typo

* add docs

* finish

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-07 15:03:58 +01:00
Pedro Cuenca
6a7f1f0965 Flax: avoid recompilation when params change (#1096)
* Do not recompile when guidance_scale changes.

* Remove debug for simplicity.

* make style

* Make guidance_scale an array.

* Make DEBUG a constant to avoid passing it down.

* Add comments for clarification.
2022-12-07 14:50:55 +01:00
Suraj Patil
170ebd288f [UNet2DConditionModel] add an option to upcast attention to fp32 (#1590)
upcast attention
2022-12-07 14:36:22 +01:00
Anton Lozhkov
dc87f526d4 Fix common tests for FP16 (#1588)
* Fix common tests for FP16

* revert
2022-12-07 14:09:51 +01:00
Fantasy-Studio
d9b5b43d46 Correct order height & width in pipeline_paint_by_example.py (#1589)
Update pipeline_paint_by_example.py
2022-12-07 13:40:56 +01:00
Anton Lozhkov
bb2d7cacc0 Add from_pretrained telemetry (#1461)
* Add from_pretrained usage logging

* Add classes

* add a telemetry notice

* macos
2022-12-07 11:56:21 +01:00
Patrick von Platen
4f3ddb6cca [Paint by Example] Better default for image width (#1587) 2022-12-07 11:43:28 +01:00
SkyTNT
4eb9ad0d1c [Community Pipeline] fix lpw_stable_diffusion (#1570)
* fix lpw_stable_diffusion

* rollback preprocess_mask resample
2022-12-07 11:20:01 +01:00
Patrick von Platen
896c98a2ae Add paint by example (#1533)
* add paint by example

* mkae loading possibel

* up

* Update src/diffusers/models/attention.py

* up

* finalize weight structure

* make example work

* make it work

* up

* up

* fix

* del

* add

* update

* Apply suggestions from code review

* correct transformer 2d

* finish

* up

* up

* up

* up

* fix

* Apply suggestions from code review

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

* Apply suggestions from code review

* up

* finish

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-07 11:06:30 +01:00
Anton Lozhkov
02d83c9ff1 Standardize fast pipeline tests with PipelineTestMixin (#1526)
* [WIP] Standardize fast pipeline tests with PipelineTestMixin

* refactor the sd tests a bit

* add more common tests

* add xformers

* add progressbar test

* cleanup

* upd fp16

* CycleDiffusionPipelineFastTests

* DanceDiffusionPipelineFastTests

* AltDiffusionPipelineFastTests

* StableDiffusion2PipelineFastTests

* StableDiffusion2InpaintPipelineFastTests

* StableDiffusionImageVariationPipelineFastTests

* StableDiffusionImg2ImgPipelineFastTests

* StableDiffusionInpaintPipelineFastTests

* remove unused mixins

* quality

* add missing inits

* try to fix mps tests

* fix mps tests

* add mps warmups

* skip for some pipelines

* style

* Update tests/test_pipelines_common.py

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-06 18:35:30 +01:00
Suraj Patil
9e1102990a [dreambooth] make collate_fn global (#1547)
make collate_fn global
2022-12-06 14:41:53 +01:00
Suraj Patil
c228331068 [examples] add check_min_version (#1550)
* add check_min_version for examples

* move __version__ to the top

* Apply suggestions from code review

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>

* fix comment

* fix error_message

* adapt the install message

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-06 14:36:50 +01:00
Patrick von Platen
ae4112d2bb Mega community pipeline (#1561)
* Mega community pipeline

* fix
2022-12-06 11:18:53 +01:00
Will Berman
af04479e85 [docs] [dreambooth training] default accelerate config (#1564) 2022-12-05 18:24:32 -08:00
Patrick von Platen
9a52e33eb6 Merge branch 'main' of https://github.com/huggingface/diffusers into main 2022-12-05 19:38:05 +00:00
Patrick von Platen
c524fd8589 correct librosa import 2022-12-05 19:37:46 +00:00
Pedro Cuenca
2cfdf37537 Fix typo (#1558)
* Fix typo in pipeline_stable_diffusion.py

Fixes a typo in a warning message

* Fix copies.

* Fix copies

Co-authored-by: Scott <scott@scottinallca.ps>
2022-12-05 20:31:35 +01:00
Patrick von Platen
62b497c418 [Docs] Correct docs (#1554) 2022-12-05 19:54:20 +01:00
Patrick von Platen
922d56a19c Correct type from int to str in conversion script sd 2022-12-05 18:51:29 +00:00
Patrick von Platen
ae854746ab [Community download] Fix cache dir (#1555)
* [Community download] Fix cache dir

* up
2022-12-05 18:52:55 +01:00
Robert Dargavel Smith
48d0123f0f add AudioDiffusionPipeline and LatentAudioDiffusionPipeline #1334 (#1426)
* add AudioDiffusionPipeline and LatentAudioDiffusionPipeline

* add docs to toc

* fix tests

* fix tests

* fix tests

* fix tests

* fix tests

* Update pr_tests.yml

Fix tests

* parent 499ff34b3e
author teticio <teticio@gmail.com> 1668765652 +0000
committer teticio <teticio@gmail.com> 1669041721 +0000

parent 499ff34b3e
author teticio <teticio@gmail.com> 1668765652 +0000
committer teticio <teticio@gmail.com> 1669041704 +0000

add colab notebook

[Flax] Fix loading scheduler from subfolder (#1319)

[FLAX] Fix loading scheduler from subfolder

Fix/Enable all schedulers for in-painting (#1331)

* inpaint fix k lms

* onnox as well

* up

Correct path to schedlure (#1322)

* [Examples] Correct path

* uP

Avoid nested fix-copies (#1332)

* Avoid nested `# Copied from` statements during `make fix-copies`

* style

Fix img2img speed with LMS-Discrete Scheduler (#896)

Casting `self.sigmas` into a different dtype (the one of original_samples) is not advisable. In my img2img pipeline this leads to a long running time in the  `integrate.quad` call later on- by long I mean more than 10x slower.

Co-authored-by: Anton Lozhkov <anton@huggingface.co>

Fix the order of casts for onnx inpainting (#1338)

Legacy Inpainting Pipeline for Onnx Models (#1237)

* Add legacy inpainting pipeline compatibility for onnx

* remove commented out line

* Add onnx legacy inpainting test

* Fix slow decorators

* pep8 styling

* isort styling

* dummy object

* ordering consistency

* style

* docstring styles

* Refactor common prompt encoding pattern

* Update tests to permanent repository home

* support all available schedulers until ONNX IO binding is available

Co-authored-by: Anton Lozhkov <aglozhkov@gmail.com>

* updated styling from PR suggested feedback

Co-authored-by: Anton Lozhkov <aglozhkov@gmail.com>

Jax infer support negative prompt (#1337)

* support negative prompts in sd jax pipeline

* pass batched neg_prompt

* only encode when negative prompt is None

Co-authored-by: Juan Acevedo <jfacevedo@google.com>

Update README.md: Minor change to Imagic code snippet, missing dir error (#1347)

Minor change to Imagic Readme

Missing dir causes an error when running the example code.

make style

change the sample model (#1352)

* Update alt_diffusion.mdx

* Update alt_diffusion.mdx

Add bit diffusion [WIP] (#971)

* Create bit_diffusion.py

Bit diffusion based on the paper, arXiv:2208.04202, Chen2022AnalogBG

* adding bit diffusion to new branch

ran tests

* tests

* tests

* tests

* tests

* removed test folders + added to README

* Update README.md

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>

* move Mel to module in pipeline construction, make librosa optional

* fix imports

* fix copy & paste error in comment

* fix style

* add missing register_to_config

* fix class docstrings

* fix class docstrings

* tweak docstrings

* tweak docstrings

* update slow test

* put trailing commas back

* respect alphabetical order

* remove LatentAudioDiffusion, make vqvae optional

* move Mel from models back to pipelines :-)

* allow loading of pretrained audiodiffusion models

* fix tests

* fix dummies

* remove reference to latent_audio_diffusion in docs

* unused import

* inherit from SchedulerMixin to make loadable

* Apply suggestions from code review

* Apply suggestions from code review

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-05 18:06:30 +01:00
Patrick von Platen
459b8ca81a Research folder (#1553)
* Research folder

* Update examples/research_projects/README.md

* up
2022-12-05 17:58:35 +01:00
Suraj Patil
bce65cd13a [refactor] make set_attention_slice recursive (#1532)
* make attn slice recursive

* remove set_attention_slice from blocks

* fix copies

* make enable_attention_slicing base class method of DiffusionPipeline

* fix set_attention_slice

* fix set_attention_slice

* fix copies

* add tests

* up

* up

* up

* update

* up

* uP

Co-authored-by: Patrick von Platen <patrick.v.platen@gmail.com>
2022-12-05 17:31:04 +01:00
Adalberto
e289998932 fix mask discrepancies in train_dreambooth_inpaint (#1529)
The mask and instance image were being cropped in different ways without --center_crop, causing the model to learn to ignore the mask in some cases. This PR fixes that and generate more consistent results.
2022-12-05 17:26:36 +01:00
Suraj Patil
634be6e53d [examples] use from_pretrained to load scheduler (#1549)
us from_pretrained to load scheduler
2022-12-05 15:32:24 +01:00
allo-
d1bcbf38ca [textual_inversion] Add an option for only saving the embeddings (#781)
[textual_inversion] Add an option to only save embeddings

Add an command line option --only_save_embeds to the example script, for
not saving the full model. Then only the learned embeddings are saved,
which can be added to the original model at runtime in a similar way as
they are created in the training script.
Saving the full model is forced when --push_to_hub is used. (Implements #759)
2022-12-05 14:45:13 +01:00
Patrick von Platen
df7cd5fe3f Update bug-report.yml 2022-12-05 14:39:35 +01:00
Naga Sai Abhinay
c28d6945b8 [Community Pipeline] Checkpoint Merger based on Automatic1111 (#1472)
* Add checkpoint_merger pipeline

* Added missing docs for a parameter.

* Fomratting fixes.

* Fixed code quality issues.

* Bug fix: Off by 1 index

* Added docs for pipeline
2022-12-05 14:36:55 +01:00
Patrick von Platen
5177e65ff0 Update bug-report.yml 2022-12-05 14:17:04 +01:00
Patrick von Platen
60ac5fc235 Update bug-report.yml 2022-12-05 14:13:02 +01:00
Patrick von Platen
19b01749f0 Update bug-report.yml 2022-12-05 14:10:25 +01:00
Patrick von Platen
a980ef2f08 Update bug-report.yml (#1548)
* Update bug-report.yml

* Update bug-report.yml

* Update bug-report.yml
2022-12-05 14:03:54 +01:00
Patrick von Platen
7932971542 [Upscaling] Fix batch size (#1525) 2022-12-05 13:28:55 +01:00
Benjamin Lefaudeux
720dbfc985 Compute embedding distances with torch.cdist (#1459)
small but mighty
2022-12-05 12:37:05 +01:00
Patrick von Platen
513fc68104 [Stable Diffusion Inpaint] Allow tensor as input image & mask (#1527)
up
2022-12-05 12:18:02 +01:00
Anton Lozhkov
cc22bda5f6 [CI] Add slow MPS tests (#1104)
* [CI] Add slow MPS tests

* fix yml

* temporarily resolve caching

* Tests: fix mps crashes.

* Skip test_load_pipeline_from_git on mps.

Not compatible with float16.

* Increase tolerance, use CPU generator, alt. slices.

* Move to nightly

* style

Co-authored-by: Pedro Cuenca <pedro@huggingface.co>
2022-12-05 11:50:24 +01:00
Ilmari Heikkinen
daebee0963 Add xformers attention to VAE (#1507)
* Add xformers attention to VAE

* Simplify VAE xformers code

* Update src/diffusers/models/attention.py

Co-authored-by: Ilmari Heikkinen <ilmari@fhtr.org>
Co-authored-by: Suraj Patil <surajp815@gmail.com>
2022-12-03 15:08:11 +01:00
Matthieu Bizien
ae368e42d2 [Proposal] Support saving to safetensors (#1494)
* Add parameter safe_serialization to DiffusionPipeline.save_pretrained

* Add option safe_serialization on ModelMixin.save_pretrained

* Add test test_save_safe_serialization

* Black

* Re-trigger the CI

* Fix doc-builder

* Validate files are saved as safetensor in test_save_safe_serialization
2022-12-02 18:33:16 +01:00
Patrick von Platen
cf4664e885 fix tests 2022-12-02 17:27:58 +00:00
Patrick von Platen
7222a8eadf make style 2022-12-02 17:18:50 +00:00
bachr
155d272cc1 Update FlaxLMSDiscreteScheduler (#1474)
- Add the missing `scale_model_input` method to `FlaxLMSDiscreteScheduler`
- Use `jnp.append` for appending to `state.derivatives`
- Use `jnp.delete` to pop from `state.derivatives`
2022-12-02 18:18:30 +01:00
Adalberto
2b30b1090f Create train_dreambooth_inpaint.py (#1091)
* Create train_dreambooth_inpaint.py

train_dreambooth.py adapted to work with the inpaint model, generating random masks during the training

* Update train_dreambooth_inpaint.py

refactored train_dreambooth_inpaint with black

* Update train_dreambooth_inpaint.py

* Update train_dreambooth_inpaint.py

* Update train_dreambooth_inpaint.py

Fix prior preservation

* add instructions to readme, fix SD2 compatibility
2022-12-02 18:06:57 +01:00
Antoine Bouthors
3ad49eeedd Fixed mask+masked_image in sd inpaint pipeline (#1516)
* Fixed mask+masked_image in sd inpaint pipeline

Those were left unset when inputs are not PIL images

* Fixed formatting
2022-12-02 17:51:51 +01:00
Patrick von Platen
769f0be8fb Finalize 2nd order schedulers (#1503)
* up

* up

* finish

* finish

* up

* up

* finish
2022-12-02 16:38:35 +01:00
Pedro Gabriel Gengo Lourenço
4f596599f4 Fix training docs to install datasets (#1476)
Fixed doc to install from training packages
2022-12-02 15:52:04 +01:00
Dhruv Naik
f57a2e0745 Fix Imagic example (#1520)
fix typo, remove incorrect arguments from .train()
2022-12-02 15:06:04 +01:00
Pedro Cuenca
3ceaa280bd Do not use torch.long in mps (#1488)
* Do not use torch.long in mps

Addresses #1056.

* Use torch.int instead of float.

* Propagate changes.

* Do not silently change float -> int.

* Propagate changes.

* Apply suggestions from code review

Co-authored-by: Anton Lozhkov <anton@huggingface.co>

Co-authored-by: Anton Lozhkov <anton@huggingface.co>
2022-12-02 13:10:17 +01:00
Benjamin Lefaudeux
a816a87a09 [refactor] Making the xformers mem-efficient attention activation recursive (#1493)
* Moving the mem efficiient attention activation to the top + recursive

* black, too bad there's no pre-commit ?

Co-authored-by: Benjamin Lefaudeux <benjamin@photoroom.com>
2022-12-02 12:30:01 +01:00
Patrick von Platen
f21415d1d9 Update conversion script to correctly handle SD 2 (#1511)
* Conversion SD 2

* finish
2022-12-02 12:28:01 +01:00
Patrick von Platen
22b9cb086b [From pretrained] Allow returning local path (#1450)
Allow returning local path
2022-12-02 12:26:39 +01:00
Will Berman
25f850a23b [docs] [dreambooth training] num_class_images clarification (#1508) 2022-12-02 12:12:28 +01:00
Will Berman
b25ae2e6ab [docs] [dreambooth training] accelerate.utils.write_basic_config (#1513) 2022-12-02 12:11:18 +01:00
Suraj Patil
0f1c24664c fix heun scheduler (#1512) 2022-12-01 22:39:57 +01:00
Anton Lozhkov
e65b71aba4 Add an explicit --image_size to the conversion script (#1509)
* Add an explicit `--image_size` to the conversion script

* style
2022-12-01 19:22:48 +01:00
Akash Gokul
a6a25ceb61 Fix Flax flip_sin_to_cos (#1369)
* Fix Flax flip_sin_to_cos

* Adding flip_sin_to_cos

Co-authored-by: Kashif Rasul <kashif.rasul@gmail.com>
2022-12-01 18:57:01 +01:00
Suraj Patil
b85bb0753e support v prediction in other schedulers (#1505)
* support v prediction in other schedulers

* v heun

* add tests for v pred

* fix tests

* fix test euler a

* v ddpm
2022-12-01 18:10:39 +01:00
fboulnois
52eb0348e5 Standardize on using image argument in all pipelines (#1361)
* feat: switch core pipelines to use image arg

* test: update tests for core pipelines

* feat: switch examples to use image arg

* docs: update docs to use image arg

* style: format code using black and doc-builder

* fix: deprecate use of init_image in all pipelines
2022-12-01 16:55:22 +01:00
Suraj Patil
2bbf8b67a7 simplyfy AttentionBlock (#1492) 2022-12-01 16:40:59 +01:00
Patrick von Platen
5a5bf7ef5a [Deprecate] Correct stacklevel (#1483)
* Correct stacklevel

* fix
2022-12-01 16:28:10 +01:00
Anton Lozhkov
9276b1e148 Replace deprecated hub utils in train_unconditional_ort (#1504)
* Replace deprecated hub utils in `train_unconditional_ort`

* typo
2022-12-01 16:00:52 +01:00
regisss
2579d42158 Add doc for Stable Diffusion on Habana Gaudi (#1496)
* Add doc for Stable Diffusion on Habana Gaudi

* Make style

* Add benchmark

* Center-align columns in the benchmark table
2022-12-01 15:43:48 +01:00
Anton Lozhkov
999044596a Bump to 0.10.0.dev0 + deprecations (#1490) 2022-11-30 15:27:56 +01:00
246 changed files with 19751 additions and 6183 deletions

View File

@@ -5,7 +5,20 @@ body:
- type: markdown
attributes:
value: |
Thanks for taking the time to fill out this bug report!
Thanks a lot for taking the time to file this issue 🤗.
Issues do not only help to improve the library, but also publicly document common problems, questions, workflows for the whole community!
Thus, issues are of the same importance as pull requests when contributing to this library ❤️.
In order to make your issue as **useful for the community as possible**, let's try to stick to some simple guidelines:
- 1. Please try to be as precise and concise as possible.
*Give your issue a fitting title. Assume that someone which very limited knowledge of diffusers can understand your issue. Add links to the source code, documentation other issues, pull requests etc...*
- 2. If your issue is about something not working, **always** provide a reproducible code snippet. The reader should be able to reproduce your issue by **only copy-pasting your code snippet into a Python shell**.
*The community cannot solve your issue if it cannot reproduce it. If your bug is related to training, add your training script and make everything needed to train public. Otherwise, just add a simple Python code snippet.*
- 3. Add the **minimum amount of code / context that is needed to understand, reproduce your issue**.
*Make the life of maintainers easy. `diffusers` is getting many issues every day. Make sure your issue is about one bug and one bug only. Make sure you add only the context, code needed to understand your issues - nothing more. Generally, every issue is a way of documenting this library, try to make it a good documentation entry.*
- type: markdown
attributes:
value: |
For more in-detail information on how to write good issues you can have a look [here](https://huggingface.co/course/chapter8/5?fw=pt)
- type: textarea
id: bug-description
attributes:
@@ -20,6 +33,8 @@ body:
label: Reproduction
description: Please provide a minimal reproducible code which we can copy/paste and reproduce the issue.
placeholder: Reproduction
validations:
required: true
- type: textarea
id: logs
attributes:

160
.github/workflows/nightly_tests.yml vendored Normal file
View File

@@ -0,0 +1,160 @@
name: Nightly tests on main
on:
schedule:
- cron: "0 0 * * *" # every day at midnight
env:
DIFFUSERS_IS_CI: yes
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
RUN_NIGHTLY: yes
jobs:
run_nightly_tests:
strategy:
fail-fast: false
matrix:
config:
- name: Nightly PyTorch CUDA tests on Ubuntu
framework: pytorch
runner: docker-gpu
image: diffusers/diffusers-pytorch-cuda
report: torch_cuda
- name: Nightly Flax TPU tests on Ubuntu
framework: flax
runner: docker-tpu
image: diffusers/diffusers-flax-tpu
report: flax_tpu
- name: Nightly ONNXRuntime CUDA tests on Ubuntu
framework: onnxruntime
runner: docker-gpu
image: diffusers/diffusers-onnxruntime-cuda
report: onnx_cuda
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ ${{ matrix.config.runner == 'docker-tpu' && '--privileged' || '--gpus 0'}}
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
if: ${{ matrix.config.runner == 'docker-gpu' }}
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m pip install -e .[quality,test]
python -m pip install -U git+https://github.com/huggingface/transformers
- name: Environment
run: |
python utils/print_env.py
- name: Run nightly PyTorch CUDA tests
if: ${{ matrix.config.framework == 'pytorch' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Run nightly Flax TPU tests
if: ${{ matrix.config.framework == 'flax' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Run nightly ONNXRuntime CUDA tests
if: ${{ matrix.config.framework == 'onnxruntime' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: ${{ matrix.config.report }}_test_reports
path: reports
run_nightly_tests_apple_m1:
name: Nightly PyTorch MPS tests on MacOS
runs-on: [ self-hosted, apple-m1 ]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Clean checkout
shell: arch -arch arm64 bash {0}
run: |
git clean -fxd
- name: Setup miniconda
uses: ./.github/actions/setup-miniconda
with:
python-version: 3.9
- name: Install dependencies
shell: arch -arch arm64 bash {0}
run: |
${CONDA_RUN} python -m pip install --upgrade pip
${CONDA_RUN} python -m pip install -e .[quality,test]
${CONDA_RUN} python -m pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
- name: Environment
shell: arch -arch arm64 bash {0}
run: |
${CONDA_RUN} python utils/print_env.py
- name: Run nightly PyTorch tests on M1 (MPS)
shell: arch -arch arm64 bash {0}
env:
HF_HOME: /System/Volumes/Data/mnt/cache
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps tests/
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_torch_mps_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_mps_test_reports
path: reports

View File

@@ -1,4 +1,4 @@
name: Run fast tests
name: Fast tests for PRs
on:
pull_request:
@@ -14,7 +14,6 @@ env:
OMP_NUM_THREADS: 4
MKL_NUM_THREADS: 4
PYTEST_TIMEOUT: 60
MPS_TORCH_VERSION: 1.13.0
jobs:
run_fast_tests:
@@ -58,8 +57,8 @@ jobs:
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate
python -m pip install -U git+https://github.com/huggingface/transformers
- name: Environment
@@ -126,8 +125,7 @@ jobs:
run: |
${CONDA_RUN} python -m pip install --upgrade pip
${CONDA_RUN} python -m pip install -e .[quality,test]
${CONDA_RUN} python -m pip install --pre torch==${MPS_TORCH_VERSION} --extra-index-url https://download.pytorch.org/whl/test/cpu
${CONDA_RUN} python -m pip install git+https://github.com/huggingface/accelerate
${CONDA_RUN} python -m pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
${CONDA_RUN} python -m pip install -U git+https://github.com/huggingface/transformers
- name: Environment
@@ -137,6 +135,9 @@ jobs:
- name: Run fast PyTorch tests on M1 (MPS)
shell: arch -arch arm64 bash {0}
env:
HF_HOME: /System/Volumes/Data/mnt/cache
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
${CONDA_RUN} python -m pytest -n 0 -s -v --make-reports=tests_torch_mps tests/

View File

@@ -1,4 +1,4 @@
name: Run all tests
name: Slow tests on main
on:
push:
@@ -10,7 +10,7 @@ env:
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 1000
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
jobs:
@@ -61,7 +61,6 @@ jobs:
- name: Install dependencies
run: |
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate
python -m pip install -U git+https://github.com/huggingface/transformers
- name: Environment
@@ -131,7 +130,6 @@ jobs:
- name: Install dependencies
run: |
python -m pip install -e .[quality,test,training]
python -m pip install git+https://github.com/huggingface/accelerate
python -m pip install -U git+https://github.com/huggingface/transformers
- name: Environment
@@ -153,4 +151,4 @@ jobs:
uses: actions/upload-artifact@v2
with:
name: examples_test_reports
path: reports
path: reports

5
.gitignore vendored
View File

@@ -165,4 +165,7 @@ tags
# DS_Store (MacOS)
.DS_Store
# RL pipelines may produce mp4 outputs
*.mp4
*.mp4
# dependencies
/transformers

View File

@@ -29,13 +29,13 @@ More precisely, 🤗 Diffusers offers:
### For PyTorch
**With `pip`**
**With `pip`** (official package)
```bash
pip install --upgrade diffusers[torch]
```
**With `conda`**
**With `conda`** (maintained by the community)
```sh
conda install -c conda-forge diffusers
@@ -79,19 +79,13 @@ In order to get started, we recommend taking a look at two notebooks:
Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [LAION](https://laion.ai/) and [RunwayML](https://runwayml.com/). It's trained on 512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/) database. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and runs on a GPU with at least 4GB VRAM.
See the [model card](https://huggingface.co/CompVis/stable-diffusion) for more information.
You need to accept the model license before downloading or using the Stable Diffusion weights. Please, visit the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5), read the license carefully and tick the checkbox if you agree. You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section](https://huggingface.co/docs/hub/security-tokens) of the documentation.
### Text-to-Image generation with Stable Diffusion
First let's install
```bash
pip install --upgrade diffusers transformers scipy
```
Run this command to log in with your HF Hub token if you haven't before (you can skip this step if you prefer to run the model locally, follow [this](#running-the-model-locally) instead)
```bash
huggingface-cli login
pip install --upgrade diffusers transformers accelerate
```
We recommend using the model in [half-precision (`fp16`)](https://pytorch.org/blog/accelerating-training-on-nvidia-gpus-with-pytorch-automatic-mixed-precision/) as it gives almost always the same results as full
@@ -101,7 +95,7 @@ precision while being roughly twice as fast and requiring half the amount of GPU
import torch
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, revision="fp16")
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
@@ -109,17 +103,16 @@ image = pipe(prompt).images[0]
```
#### Running the model locally
If you don't want to login to Hugging Face, you can also simply download the model folder
(after having [accepted the license](https://huggingface.co/runwayml/stable-diffusion-v1-5)) and pass
the path to the local folder to the `StableDiffusionPipeline`.
You can also simply download the model folder and pass the path to the local folder to the `StableDiffusionPipeline`.
```
git lfs install
git clone https://huggingface.co/runwayml/stable-diffusion-v1-5
```
Assuming the folder is stored locally under `./stable-diffusion-v1-5`, you can also run stable diffusion
without requiring an authentication token:
Assuming the folder is stored locally under `./stable-diffusion-v1-5`, you can run stable diffusion
as follows:
```python
pipe = StableDiffusionPipeline.from_pretrained("./stable-diffusion-v1-5")
@@ -134,11 +127,7 @@ to using `fp16`.
The following snippet should result in less than 4GB VRAM.
```python
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
@@ -164,7 +153,6 @@ If you want to run Stable Diffusion on CPU or you want to have maximum precision
please run the model in the default *full-precision* setting:
```python
# make sure you're logged in with `huggingface-cli login`
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
@@ -262,11 +250,8 @@ from diffusers import StableDiffusionImg2ImgPipeline
# load the pipeline
device = "cuda"
model_id_or_path = "runwayml/stable-diffusion-v1-5"
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
model_id_or_path,
revision="fp16",
torch_dtype=torch.float16,
)
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(model_id_or_path, torch_dtype=torch.float16)
# or download via git clone https://huggingface.co/runwayml/stable-diffusion-v1-5
# and pass `model_id_or_path="./stable-diffusion-v1-5"`.
pipe = pipe.to(device)
@@ -280,7 +265,7 @@ init_image = init_image.resize((768, 512))
prompt = "A fantasy landscape, trending on artstation"
images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
images[0].save("fantasy_landscape.png")
```
@@ -288,10 +273,7 @@ You can also run this example on colab [![Open In Colab](https://colab.research.
### In-painting using Stable Diffusion
The `StableDiffusionInpaintPipeline` lets you edit specific parts of an image by providing a mask and a text prompt. It uses a model optimized for this particular task, whose license you need to accept before use.
Please, visit the [model card](https://huggingface.co/runwayml/stable-diffusion-inpainting), read the license carefully and tick the checkbox if you agree. Note that this is an additional license, you need to accept it even if you accepted the text-to-image Stable Diffusion license in the past. You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section](https://huggingface.co/docs/hub/security-tokens) of the documentation.
The `StableDiffusionInpaintPipeline` lets you edit specific parts of an image by providing a mask and a text prompt.
```python
import PIL
@@ -311,11 +293,7 @@ mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data
init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
pipe = StableDiffusionInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = StableDiffusionInpaintPipeline.from_pretrained("runwayml/stable-diffusion-inpainting", torch_dtype=torch.float16)
pipe = pipe.to("cuda")
prompt = "Face of a yellow cat, high resolution, sitting on a park bench"
@@ -324,11 +302,8 @@ image = pipe(prompt=prompt, image=init_image, mask_image=mask_image).images[0]
### Tweak prompts reusing seeds and latents
You can generate your own latents to reproduce results, or tweak your prompt on a specific result you liked. [This notebook](https://github.com/pcuenca/diffusers-examples/blob/main/notebooks/stable-diffusion-seeds.ipynb) shows how to do it step by step. You can also run it in Google Colab [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/pcuenca/diffusers-examples/blob/main/notebooks/stable-diffusion-seeds.ipynb).
For more details, check out [the Stable Diffusion notebook](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb) [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/stable_diffusion.ipynb)
and have a look into the [release notes](https://github.com/huggingface/diffusers/releases/tag/v0.2.0).
You can generate your own latents to reproduce results, or tweak your prompt on a specific result you liked.
Please have a look at [Reusing seeds for deterministic generation](https://huggingface.co/docs/diffusers/main/en/using-diffusers/reusing_seeds).
## Fine-Tuning Stable Diffusion

View File

@@ -11,6 +11,7 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -33,6 +34,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
librosa \
modelcards \
numpy \
scipy \

View File

@@ -11,6 +11,7 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -35,6 +36,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
librosa \
modelcards \
numpy \
scipy \

View File

@@ -11,6 +11,7 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -33,6 +34,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
librosa \
modelcards \
numpy \
scipy \

View File

@@ -11,6 +11,7 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -33,6 +34,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
librosa \
modelcards \
numpy \
scipy \

View File

@@ -11,6 +11,7 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -32,6 +33,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
librosa \
modelcards \
numpy \
scipy \

View File

@@ -11,6 +11,7 @@ RUN apt update && \
git-lfs \
curl \
ca-certificates \
libsndfile1-dev \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -32,6 +33,7 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip && \
datasets \
hf-doc-builder \
huggingface-hub \
librosa \
modelcards \
numpy \
scipy \

266
docs/README.md Normal file
View File

@@ -0,0 +1,266 @@
<!---
Copyright 2022- The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License");
you may not use this file except in compliance with the License.
You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software
distributed under the License is distributed on an "AS IS" BASIS,
WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
See the License for the specific language governing permissions and
limitations under the License.
-->
# Generating the documentation
To generate the documentation, you first have to build it. Several packages are necessary to build the doc,
you can install them with the following command, at the root of the code repository:
```bash
pip install -e ".[docs]"
```
Then you need to install our open source documentation builder tool:
```bash
pip install git+https://github.com/huggingface/doc-builder
```
---
**NOTE**
You only need to generate the documentation to inspect it locally (if you're planning changes and want to
check how they look before committing for instance). You don't have to commit the built documentation.
---
## Previewing the documentation
To preview the docs, first install the `watchdog` module with:
```bash
pip install watchdog
```
Then run the following command:
```bash
doc-builder preview {package_name} {path_to_docs}
```
For example:
```bash
doc-builder preview diffusers docs/source/
```
The docs will be viewable at [http://localhost:3000](http://localhost:3000). You can also preview the docs once you have opened a PR. You will see a bot add a comment to a link where the documentation with your changes lives.
---
**NOTE**
The `preview` command only works with existing doc files. When you add a completely new file, you need to update `_toctree.yml` & restart `preview` command (`ctrl-c` to stop it & call `doc-builder preview ...` again).
---
## Adding a new element to the navigation bar
Accepted files are Markdown (.md or .mdx).
Create a file with its extension and put it in the source directory. You can then link it to the toc-tree by putting
the filename without the extension in the [`_toctree.yml`](https://github.com/huggingface/diffusers/blob/main/docs/source/_toctree.yml) file.
## Renaming section headers and moving sections
It helps to keep the old links working when renaming the section header and/or moving sections from one document to another. This is because the old links are likely to be used in Issues, Forums, and Social media and it'd make for a much more superior user experience if users reading those months later could still easily navigate to the originally intended information.
Therefore, we simply keep a little map of moved sections at the end of the document where the original section was. The key is to preserve the original anchor.
So if you renamed a section from: "Section A" to "Section B", then you can add at the end of the file:
```
Sections that were moved:
[ <a href="#section-b">Section A</a><a id="section-a"></a> ]
```
and of course, if you moved it to another file, then:
```
Sections that were moved:
[ <a href="../new-file#section-b">Section A</a><a id="section-a"></a> ]
```
Use the relative style to link to the new file so that the versioned docs continue to work.
For an example of a rich moved section set please see the very end of [the transformers Trainer doc](https://github.com/huggingface/transformers/blob/main/docs/source/en/main_classes/trainer.mdx).
## Writing Documentation - Specification
The `huggingface/diffusers` documentation follows the
[Google documentation](https://sphinxcontrib-napoleon.readthedocs.io/en/latest/example_google.html) style for docstrings,
although we can write them directly in Markdown.
### Adding a new tutorial
Adding a new tutorial or section is done in two steps:
- Add a new file under `docs/source`. This file can either be ReStructuredText (.rst) or Markdown (.md).
- Link that file in `docs/source/_toctree.yml` on the correct toc-tree.
Make sure to put your new file under the proper section. It's unlikely to go in the first section (*Get Started*), so
depending on the intended targets (beginners, more advanced users, or researchers) it should go in sections two, three, or four.
### Adding a new pipeline/scheduler
When adding a new pipeline:
- create a file `xxx.mdx` under `docs/source/api/pipelines` (don't hesitate to copy an existing file as template).
- Link that file in (*Diffusers Summary*) section in `docs/source/api/pipelines/overview.mdx`, along with the link to the paper, and a colab notebook (if available).
- Write a short overview of the diffusion model:
- Overview with paper & authors
- Paper abstract
- Tips and tricks and how to use it best
- Possible an end-to-end example of how to use it
- Add all the pipeline classes that should be linked in the diffusion model. These classes should be added using our Markdown syntax. Usually as follows:
```
## XXXPipeline
[[autodoc]] XXXPipeline
```
This will include every public method of the pipeline that is documented. You can specify which methods should be in the docs:
```
## XXXPipeline
[[autodoc]] XXXPipeline
- __call__
```
You can follow the same process to create a new scheduler under the `docs/source/api/schedulers` folder
### Writing source documentation
Values that should be put in `code` should either be surrounded by backticks: \`like so\`. Note that argument names
and objects like True, None, or any strings should usually be put in `code`.
When mentioning a class, function, or method, it is recommended to use our syntax for internal links so that our tool
adds a link to its documentation with this syntax: \[\`XXXClass\`\] or \[\`function\`\]. This requires the class or
function to be in the main package.
If you want to create a link to some internal class or function, you need to
provide its path. For instance: \[\`pipeline_utils.ImagePipelineOutput\`\]. This will be converted into a link with
`pipeline_utils.ImagePipelineOutput` in the description. To get rid of the path and only keep the name of the object you are
linking to in the description, add a ~: \[\`~pipeline_utils.ImagePipelineOutput\`\] will generate a link with `ImagePipelineOutput` in the description.
The same works for methods so you can either use \[\`XXXClass.method\`\] or \[~\`XXXClass.method\`\].
#### Defining arguments in a method
Arguments should be defined with the `Args:` (or `Arguments:` or `Parameters:`) prefix, followed by a line return and
an indentation. The argument should be followed by its type, with its shape if it is a tensor, a colon, and its
description:
```
Args:
n_layers (`int`): The number of layers of the model.
```
If the description is too long to fit in one line, another indentation is necessary before writing the description
after the argument.
Here's an example showcasing everything so far:
```
Args:
input_ids (`torch.LongTensor` of shape `(batch_size, sequence_length)`):
Indices of input sequence tokens in the vocabulary.
Indices can be obtained using [`AlbertTokenizer`]. See [`~PreTrainedTokenizer.encode`] and
[`~PreTrainedTokenizer.__call__`] for details.
[What are input IDs?](../glossary#input-ids)
```
For optional arguments or arguments with defaults we follow the following syntax: imagine we have a function with the
following signature:
```
def my_function(x: str = None, a: float = 1):
```
then its documentation should look like this:
```
Args:
x (`str`, *optional*):
This argument controls ...
a (`float`, *optional*, defaults to 1):
This argument is used to ...
```
Note that we always omit the "defaults to \`None\`" when None is the default for any argument. Also note that even
if the first line describing your argument type and its default gets long, you can't break it on several lines. You can
however write as many lines as you want in the indented description (see the example above with `input_ids`).
#### Writing a multi-line code block
Multi-line code blocks can be useful for displaying examples. They are done between two lines of three backticks as usual in Markdown:
````
```
# first line of code
# second line
# etc
```
````
#### Writing a return block
The return block should be introduced with the `Returns:` prefix, followed by a line return and an indentation.
The first line should be the type of the return, followed by a line return. No need to indent further for the elements
building the return.
Here's an example of a single value return:
```
Returns:
`List[int]`: A list of integers in the range [0, 1] --- 1 for a special token, 0 for a sequence token.
```
Here's an example of a tuple return, comprising several objects:
```
Returns:
`tuple(torch.FloatTensor)` comprising various elements depending on the configuration ([`BertConfig`]) and inputs:
- ** loss** (*optional*, returned when `masked_lm_labels` is provided) `torch.FloatTensor` of shape `(1,)` --
Total loss is the sum of the masked language modeling loss and the next sequence prediction (classification) loss.
- **prediction_scores** (`torch.FloatTensor` of shape `(batch_size, sequence_length, config.vocab_size)`) --
Prediction scores of the language modeling head (scores for each vocabulary token before SoftMax).
```
#### Adding an image
Due to the rapidly growing repository, it is important to make sure that no files that would significantly weigh down the repository are added. This includes images, videos, and other non-text files. We prefer to leverage a hf.co hosted `dataset` like
the ones hosted on [`hf-internal-testing`](https://huggingface.co/hf-internal-testing) in which to place these files and reference
them by URL. We recommend putting them in the following dataset: [huggingface/documentation-images](https://huggingface.co/datasets/huggingface/documentation-images).
If an external contribution, feel free to add the images to your PR and ask a Hugging Face member to migrate your images
to this dataset.
## Styling the docstring
We have an automatic script running with the `make style` command that will make sure that:
- the docstrings fully take advantage of the line width
- all code examples are formatted using black, like the code of the Transformers library
This script may have some weird failures if you made a syntax mistake or if you uncover a bug. Therefore, it's
recommended to commit your changes before running `make style`, so you can revert the changes done by that script
easily.

View File

@@ -26,6 +26,10 @@
title: "Text-Guided Image-to-Image"
- local: using-diffusers/inpaint
title: "Text-Guided Image-Inpainting"
- local: using-diffusers/depth2img
title: "Text-Guided Depth-to-Image"
- local: using-diffusers/reusing_seeds
title: "Reusing seeds for deterministic generation"
- local: using-diffusers/custom_pipeline_examples
title: "Community Pipelines"
- local: using-diffusers/contribute_pipeline
@@ -43,12 +47,16 @@
- sections:
- local: optimization/fp16
title: "Memory and Speed"
- local: optimization/xformers
title: "xFormers"
- local: optimization/onnx
title: "ONNX"
- local: optimization/open_vino
title: "OpenVINO"
- local: optimization/mps
title: "MPS"
- local: optimization/habana
title: "Habana Gaudi"
title: "Optimization/Special Hardware"
- sections:
- local: training/overview
@@ -74,8 +82,6 @@
- sections:
- local: api/models
title: "Models"
- local: api/schedulers
title: "Schedulers"
- local: api/diffusion_pipeline
title: "Diffusion Pipeline"
- local: api/logging
@@ -100,6 +106,8 @@
title: "Latent Diffusion"
- local: api/pipelines/latent_diffusion_uncond
title: "Unconditional Latent Diffusion"
- local: api/pipelines/paint_by_example
title: "PaintByExample"
- local: api/pipelines/pndm
title: "PNDM"
- local: api/pipelines/score_sde_ve
@@ -114,13 +122,55 @@
title: "Stochastic Karras VE"
- local: api/pipelines/dance_diffusion
title: "Dance Diffusion"
- local: api/pipelines/unclip
title: "UnCLIP"
- local: api/pipelines/versatile_diffusion
title: "Versatile Diffusion"
- local: api/pipelines/vq_diffusion
title: "VQ Diffusion"
- local: api/pipelines/repaint
title: "RePaint"
- local: api/pipelines/audio_diffusion
title: "Audio Diffusion"
title: "Pipelines"
- sections:
- local: api/schedulers/overview
title: "Overview"
- local: api/schedulers/ddim
title: "DDIM"
- local: api/schedulers/ddpm
title: "DDPM"
- local: api/schedulers/singlestep_dpm_solver
title: "Singlestep DPM-Solver"
- local: api/schedulers/multistep_dpm_solver
title: "Multistep DPM-Solver"
- local: api/schedulers/heun
title: "Heun Scheduler"
- local: api/schedulers/dpm_discrete
title: "DPM Discrete Scheduler"
- local: api/schedulers/dpm_discrete_ancestral
title: "DPM Discrete Scheduler with ancestral sampling"
- local: api/schedulers/stochastic_karras_ve
title: "Stochastic Kerras VE"
- local: api/schedulers/lms_discrete
title: "Linear Multistep"
- local: api/schedulers/pndm
title: "PNDM"
- local: api/schedulers/score_sde_ve
title: "VE-SDE"
- local: api/schedulers/ipndm
title: "IPNDM"
- local: api/schedulers/score_sde_vp
title: "VP-SDE"
- local: api/schedulers/euler
title: "Euler scheduler"
- local: api/schedulers/euler_ancestral
title: "Euler Ancestral Scheduler"
- local: api/schedulers/vq_diffusion
title: "VQDiffusionScheduler"
- local: api/schedulers/repaint
title: "RePaint Scheduler"
title: "Schedulers"
- sections:
- local: api/experimental/rl
title: "RL Planning"

View File

@@ -58,6 +58,12 @@ The models are built on the base class ['ModelMixin'] that is a `torch.nn.module
## Transformer2DModelOutput
[[autodoc]] models.attention.Transformer2DModelOutput
## PriorTransformer
[[autodoc]] models.prior_transformer.PriorTransformer
## PriorTransformerOutput
[[autodoc]] models.prior_transformer.PriorTransformerOutput
## FlaxModelMixin
[[autodoc]] FlaxModelMixin

View File

@@ -0,0 +1,102 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Audio Diffusion
## Overview
[Audio Diffusion](https://github.com/teticio/audio-diffusion) by Robert Dargavel Smith.
Audio Diffusion leverages the recent advances in image generation using diffusion models by converting audio samples to
and from mel spectrogram images.
The original codebase of this implementation can be found [here](https://github.com/teticio/audio-diffusion), including
training scripts and example notebooks.
## Available Pipelines:
| Pipeline | Tasks | Colab
|---|---|:---:|
| [pipeline_audio_diffusion.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/audio_diffusion/pipeline_audio_diffusion.py) | *Unconditional Audio Generation* | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/teticio/audio-diffusion/blob/master/notebooks/audio_diffusion_pipeline.ipynb) |
## Examples:
### Audio Diffusion
```python
import torch
from IPython.display import Audio
from diffusers import DiffusionPipeline
device = "cuda" if torch.cuda.is_available() else "cpu"
pipe = DiffusionPipeline.from_pretrained("teticio/audio-diffusion-256").to(device)
output = pipe()
display(output.images[0])
display(Audio(output.audios[0], rate=mel.get_sample_rate()))
```
### Latent Audio Diffusion
```python
import torch
from IPython.display import Audio
from diffusers import DiffusionPipeline
device = "cuda" if torch.cuda.is_available() else "cpu"
pipe = DiffusionPipeline.from_pretrained("teticio/latent-audio-diffusion-256").to(device)
output = pipe()
display(output.images[0])
display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
```
### Audio Diffusion with DDIM (faster)
```python
import torch
from IPython.display import Audio
from diffusers import DiffusionPipeline
device = "cuda" if torch.cuda.is_available() else "cpu"
pipe = DiffusionPipeline.from_pretrained("teticio/audio-diffusion-ddim-256").to(device)
output = pipe()
display(output.images[0])
display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
```
### Variations, in-painting, out-painting etc.
```python
output = pipe(
raw_audio=output.audios[0, 0],
start_step=int(pipe.get_default_steps() / 2),
mask_start_secs=1,
mask_end_secs=1,
)
display(output.images[0])
display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))
```
## AudioDiffusionPipeline
[[autodoc]] AudioDiffusionPipeline
- __call__
- encode
- slerp
## Mel
[[autodoc]] Mel
- audio_slice_to_image
- image_to_audio

View File

@@ -57,7 +57,7 @@ prompt = "An astronaut riding an elephant"
image = pipe(
prompt=prompt,
source_prompt=source_prompt,
init_image=init_image,
image=init_image,
num_inference_steps=100,
eta=0.1,
strength=0.8,
@@ -83,7 +83,7 @@ torch.manual_seed(0)
image = pipe(
prompt=prompt,
source_prompt=source_prompt,
init_image=init_image,
image=init_image,
num_inference_steps=100,
eta=0.1,
strength=0.85,

View File

@@ -44,29 +44,32 @@ available a colab notebook to directly try them out.
| Pipeline | Paper | Tasks | Colab
|---|---|:---:|:---:|
| [alt_diffusion](./api/pipelines/alt_diffusion) | [**AltDiffusion**](https://arxiv.org/abs/2211.06679) | Image-to-Image Text-Guided Generation | -
| [cycle_diffusion](./api/pipelines/cycle_diffusion) | [**Cycle Diffusion**](https://arxiv.org/abs/2210.05559) | Image-to-Image Text-Guided Generation |
| [dance_diffusion](./api/pipelines/dance_diffusion) | [**Dance Diffusion**](https://github.com/williamberman/diffusers.git) | Unconditional Audio Generation |
| [ddpm](./api/pipelines/ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) | Unconditional Image Generation |
| [ddim](./api/pipelines/ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) | Unconditional Image Generation |
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation |
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image |
| [latent_diffusion_uncond](./api/pipelines/latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation |
| [pndm](./api/pipelines/pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation |
| [score_sde_ve](./api/pipelines/score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [score_sde_vp](./api/pipelines/score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
| [stable_diffusion](./api/pipelines/stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-to-Image Generation |
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting |
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
| [stable_diffusion_safe](./api/pipelines/stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
| [stochastic_karras_ve](./api/pipelines/stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |
| [vq_diffusion](./api/pipelines/vq_diffusion) | [Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation |
| [alt_diffusion](./alt_diffusion) | [**AltDiffusion**](https://arxiv.org/abs/2211.06679) | Image-to-Image Text-Guided Generation | -
| [audio_diffusion](./audio_diffusion) | [**Audio Diffusion**](https://github.com/teticio/audio_diffusion.git) | Unconditional Audio Generation |
| [cycle_diffusion](./cycle_diffusion) | [**Cycle Diffusion**](https://arxiv.org/abs/2210.05559) | Image-to-Image Text-Guided Generation |
| [dance_diffusion](./dance_diffusion) | [**Dance Diffusion**](https://github.com/williamberman/diffusers.git) | Unconditional Audio Generation |
| [ddpm](./ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) | Unconditional Image Generation |
| [ddim](./ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) | Unconditional Image Generation |
| [latent_diffusion](./latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation |
| [latent_diffusion](./latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image |
| [latent_diffusion_uncond](./latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation |
| [paint_by_example](./paint_by_example) | [**Paint by Example: Exemplar-based Image Editing with Diffusion Models**](https://arxiv.org/abs/2211.13227) | Image-Guided Image Inpainting |
| [pndm](./pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation |
| [score_sde_ve](./score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [score_sde_vp](./score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [stable_diffusion](./stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-to-Image Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/training_example.ipynb)
| [stable_diffusion](./stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Image-to-Image Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/image_2_image_using_diffusers.ipynb)
| [stable_diffusion](./stable_diffusion) | [**Stable Diffusion**](https://stability.ai/blog/stable-diffusion-public-release) | Text-Guided Image Inpainting | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/in_painting_with_stable_diffusion_using_diffusers.ipynb)
| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-to-Image Generation |
| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Image Inpainting |
| [stable_diffusion_2](./stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
| [stable_diffusion_safe](./stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
| [stochastic_karras_ve](./stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
| [unclip](./unclip) | [Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) | Text-to-Image Generation |
| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation |
| [versatile_diffusion](./versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |
| [vq_diffusion](./vq_diffusion) | [Vector Quantized Diffusion Model for Text-to-Image Synthesis](https://arxiv.org/abs/2111.14822) | Text-to-Image Generation |
**Note**: Pipelines are simple examples of how to play around with the diffusion systems as described in the corresponding papers.
@@ -136,9 +139,9 @@ from diffusers import StableDiffusionImg2ImgPipeline
# load the pipeline
device = "cuda"
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", revision="fp16", torch_dtype=torch.float16
).to(device)
pipe = StableDiffusionImg2ImgPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16).to(
device
)
# let's download an initial image
url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
@@ -149,7 +152,7 @@ init_image = init_image.resize((768, 512))
prompt = "A fantasy landscape, trending on artstation"
images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
images[0].save("fantasy_landscape.png")
```
@@ -186,7 +189,6 @@ mask_image = download_image(mask_url).resize((512, 512))
pipe = StableDiffusionInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")

View File

@@ -0,0 +1,73 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# PaintByExample
## Overview
[Paint by Example: Exemplar-based Image Editing with Diffusion Models](https://arxiv.org/abs/2211.13227) by Binxin Yang, Shuyang Gu, Bo Zhang, Ting Zhang, Xuejin Chen, Xiaoyan Sun, Dong Chen, Fang Wen
The abstract of the paper is the following:
*Language-guided image editing has achieved great success recently. In this paper, for the first time, we investigate exemplar-guided image editing for more precise control. We achieve this goal by leveraging self-supervised training to disentangle and re-organize the source image and the exemplar. However, the naive approach will cause obvious fusing artifacts. We carefully analyze it and propose an information bottleneck and strong augmentations to avoid the trivial solution of directly copying and pasting the exemplar image. Meanwhile, to ensure the controllability of the editing process, we design an arbitrary shape mask for the exemplar image and leverage the classifier-free guidance to increase the similarity to the exemplar image. The whole framework involves a single forward of the diffusion model without any iterative optimization. We demonstrate that our method achieves an impressive performance and enables controllable editing on in-the-wild images with high fidelity.*
The original codebase can be found [here](https://github.com/Fantasy-Studio/Paint-by-Example).
## Available Pipelines:
| Pipeline | Tasks | Colab
|---|---|:---:|
| [pipeline_paint_by_example.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/paint_by_example/pipeline_paint_by_example.py) | *Image-Guided Image Painting* | - |
## Tips
- PaintByExample is supported by the official [Fantasy-Studio/Paint-by-Example](https://huggingface.co/Fantasy-Studio/Paint-by-Example) checkpoint. The checkpoint has been warm-started from the [CompVis/stable-diffusion-v1-4](https://huggingface.co/CompVis/stable-diffusion-v1-4) and with the objective to inpaint partly masked images conditioned on example / reference images
- To quickly demo *PaintByExample*, please have a look at [this demo](https://huggingface.co/spaces/Fantasy-Studio/Paint-by-Example)
- You can run the following code snippet as an example:
```python
# !pip install diffusers transformers
import PIL
import requests
import torch
from io import BytesIO
from diffusers import DiffusionPipeline
def download_image(url):
response = requests.get(url)
return PIL.Image.open(BytesIO(response.content)).convert("RGB")
img_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/image/example_1.png"
mask_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/mask/example_1.png"
example_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/reference/example_1.jpg"
init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
example_image = download_image(example_url).resize((512, 512))
pipe = DiffusionPipeline.from_pretrained(
"Fantasy-Studio/Paint-by-Example",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
image = pipe(image=init_image, mask_image=mask_image, example_image=example_image).images[0]
image
```
## PaintByExamplePipeline
[[autodoc]] pipelines.paint_by_example.pipeline_paint_by_example.PaintByExamplePipeline
- __call__

View File

@@ -97,6 +97,14 @@ If you want to use all possible use cases in a single `DiffusionPipeline` you ca
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
## StableDiffusionDepth2ImgPipeline
[[autodoc]] StableDiffusionDepth2ImgPipeline
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
## StableDiffusionImageVariationPipeline
[[autodoc]] StableDiffusionImageVariationPipeline
- __call__

View File

@@ -30,6 +30,7 @@ Note that the architecture is more or less identical to [Stable Diffusion 1](./a
- *Text-to-Image (768x768 resolution)*: [stabilityai/stable-diffusion-2](https://huggingface.co/stabilityai/stable-diffusion-2) with [`StableDiffusionPipeline`]
- *Image Inpainting (512x512 resolution)*: [stabilityai/stable-diffusion-2-inpainting](https://huggingface.co/stabilityai/stable-diffusion-2-inpainting) with [`StableDiffusionInpaintPipeline`]
- *Image Upscaling (x4 resolution resolution)*: [stable-diffusion-x4-upscaler](https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler) [`StableDiffusionUpscalePipeline`]
- *Depth-to-Image (512x512 resolution)*: [stabilityai/stable-diffusion-2-depth](https://huggingface.co/stabilityai/stable-diffusion-2-depth) with [`StableDiffusionDepth2ImagePipeline`]
We recommend using the [`DPMSolverMultistepScheduler`] as it's currently the fastest scheduler there is.
@@ -112,7 +113,7 @@ import torch
# load model and scheduler
model_id = "stabilityai/stable-diffusion-x4-upscaler"
pipeline = StableDiffusionUpscalePipeline.from_pretrained(model_id, revision="fp16", torch_dtype=torch.float16)
pipeline = StableDiffusionUpscalePipeline.from_pretrained(model_id, torch_dtype=torch.float16)
pipeline = pipeline.to("cuda")
# let's download an image
@@ -125,6 +126,37 @@ upscaled_image = pipeline(prompt=prompt, image=low_res_img).images[0]
upscaled_image.save("upsampled_cat.png")
```
- *Depth-Guided Text-to-Image*: [stabilityai/stable-diffusion-2-depth](https://huggingface.co/stabilityai/stable-diffusion-2-depth) [`StableDiffusionDepth2ImagePipeline`]
**Installation**
```bash
!pip install -U git+https://github.com/huggingface/transformers.git
!pip install diffusers[torch]
```
**Example**
```python
import torch
import requests
from PIL import Image
from diffusers import StableDiffusionDepth2ImgPipeline
pipe = StableDiffusionDepth2ImgPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-depth",
torch_dtype=torch.float16,
).to("cuda")
url = "http://images.cocodataset.org/val2017/000000039769.jpg"
init_image = Image.open(requests.get(url, stream=True).raw)
prompt = "two tigers"
n_propmt = "bad, deformed, ugly, bad anotomy"
image = pipe(prompt=prompt, image=init_image, negative_prompt=n_propmt, strength=0.7).images[0]
```
### How to load and use different schedulers.
The stable diffusion pipeline uses [`DDIMScheduler`] scheduler by default. But `diffusers` provides many other schedulers that can be used with the stable diffusion pipeline such as [`PNDMScheduler`], [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`] etc.

View File

@@ -0,0 +1,31 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# unCLIP
## Overview
[Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) by Aditya Ramesh, Prafulla Dhariwal, Alex Nichol, Casey Chu, Mark Chen
The abstract of the paper is the following:
Contrastive models like CLIP have been shown to learn robust representations of images that capture both semantics and style. To leverage these representations for image generation, we propose a two-stage model: a prior that generates a CLIP image embedding given a text caption, and a decoder that generates an image conditioned on the image embedding. We show that explicitly generating image representations improves image diversity with minimal loss in photorealism and caption similarity. Our decoders conditioned on image representations can also produce variations of an image that preserve both its semantics and style, while varying the non-essential details absent from the image representation. Moreover, the joint embedding space of CLIP enables language-guided image manipulations in a zero-shot fashion. We use diffusion models for the decoder and experiment with both autoregressive and diffusion models for the prior, finding that the latter are computationally more efficient and produce higher-quality samples.
The unCLIP model in diffusers comes from kakaobrain's karlo and the original codebase can be found [here](https://github.com/kakaobrain/karlo). Additionally, lucidrains has a DALL-E 2 recreation [here](https://github.com/lucidrains/DALLE2-pytorch).
## Available Pipelines:
| Pipeline | Tasks | Colab
|---|---|:---:|
| [pipeline_unclip.py](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/unclip/pipeline_unclip.py) | *Text-to-Image Generation* | - |
## UnCLIPPipeline
[[autodoc]] pipelines.unclip.pipeline_unclip.UnCLIPPipeline
- __call__

View File

@@ -1,151 +0,0 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Schedulers
Diffusers contains multiple pre-built schedule functions for the diffusion process.
## What is a scheduler?
The schedule functions, denoted *Schedulers* in the library take in the output of a trained model, a sample which the diffusion process is iterating on, and a timestep to return a denoised sample. That's why schedulers may also be called *Samplers* in other diffusion models implementations.
- Schedulers define the methodology for iteratively adding noise to an image or for updating a sample based on model outputs.
- adding noise in different manners represent the algorithmic processes to train a diffusion model by adding noise to images.
- for inference, the scheduler defines how to update a sample based on an output from a pretrained model.
- Schedulers are often defined by a *noise schedule* and an *update rule* to solve the differential equation solution.
### Discrete versus continuous schedulers
All schedulers take in a timestep to predict the updated version of the sample being diffused.
The timesteps dictate where in the diffusion process the step is, where data is generated by iterating forward in time and inference is executed by propagating backwards through timesteps.
Different algorithms use timesteps that both discrete (accepting `int` inputs), such as the [`DDPMScheduler`] or [`PNDMScheduler`], and continuous (accepting `float` inputs), such as the score-based schedulers [`ScoreSdeVeScheduler`] or [`ScoreSdeVpScheduler`].
## Designing Re-usable schedulers
The core design principle between the schedule functions is to be model, system, and framework independent.
This allows for rapid experimentation and cleaner abstractions in the code, where the model prediction is separated from the sample update.
To this end, the design of schedulers is such that:
- Schedulers can be used interchangeably between diffusion models in inference to find the preferred trade-off between speed and generation quality.
- Schedulers are currently by default in PyTorch, but are designed to be framework independent (partial Jax support currently exists).
## API
The core API for any new scheduler must follow a limited structure.
- Schedulers should provide one or more `def step(...)` functions that should be called to update the generated sample iteratively.
- Schedulers should provide a `set_timesteps(...)` method that configures the parameters of a schedule function for a specific inference task.
- Schedulers should be framework-specific.
The base class [`SchedulerMixin`] implements low level utilities used by multiple schedulers.
### SchedulerMixin
[[autodoc]] SchedulerMixin
### SchedulerOutput
The class [`SchedulerOutput`] contains the outputs from any schedulers `step(...)` call.
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput
### Implemented Schedulers
#### Denoising diffusion implicit models (DDIM)
Original paper can be found here.
[[autodoc]] DDIMScheduler
#### Denoising diffusion probabilistic models (DDPM)
Original paper can be found [here](https://arxiv.org/abs/2010.02502).
[[autodoc]] DDPMScheduler
#### Multistep DPM-Solver
Original paper can be found [here](https://arxiv.org/abs/2206.00927) and the [improved version](https://arxiv.org/abs/2211.01095). The original implementation can be found [here](https://github.com/LuChengTHU/dpm-solver).
[[autodoc]] DPMSolverMultistepScheduler
#### Variance exploding, stochastic sampling from Karras et. al
Original paper can be found [here](https://arxiv.org/abs/2006.11239).
[[autodoc]] KarrasVeScheduler
#### Linear multistep scheduler for discrete beta schedules
Original implementation can be found [here](https://arxiv.org/abs/2206.00364).
[[autodoc]] LMSDiscreteScheduler
#### Pseudo numerical methods for diffusion models (PNDM)
Original implementation can be found [here](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181).
[[autodoc]] PNDMScheduler
#### variance exploding stochastic differential equation (VE-SDE) scheduler
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
[[autodoc]] ScoreSdeVeScheduler
#### improved pseudo numerical methods for diffusion models (iPNDM)
Original implementation can be found [here](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296).
[[autodoc]] IPNDMScheduler
#### variance preserving stochastic differential equation (VP-SDE) scheduler
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
<Tip warning={true}>
Score SDE-VP is under construction.
</Tip>
[[autodoc]] schedulers.scheduling_sde_vp.ScoreSdeVpScheduler
#### Euler scheduler
Euler scheduler (Algorithm 2) from the paper [Elucidating the Design Space of Diffusion-Based Generative Models](https://arxiv.org/abs/2206.00364) by Karras et al. (2022). Based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L51) implementation by Katherine Crowson.
Fast scheduler which often times generates good outputs with 20-30 steps.
[[autodoc]] EulerDiscreteScheduler
#### Euler Ancestral scheduler
Ancestral sampling with Euler method steps. Based on the original (k-diffusion)[https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72] implementation by Katherine Crowson.
Fast scheduler which often times generates good outputs with 20-30 steps.
[[autodoc]] EulerAncestralDiscreteScheduler
#### VQDiffusionScheduler
Original paper can be found [here](https://arxiv.org/abs/2111.14822)
[[autodoc]] VQDiffusionScheduler
#### RePaint scheduler
DDPM-based inpainting scheduler for unsupervised inpainting with extreme masks.
Intended for use with [`RePaintPipeline`].
Based on the paper [RePaint: Inpainting using Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2201.09865)
and the original implementation by Andreas Lugmayr et al.: https://github.com/andreas128/RePaint
[[autodoc]] RePaintScheduler

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# Denoising diffusion implicit models (DDIM)
## Overview
[Denoising Diffusion Implicit Models](https://arxiv.org/abs/2010.02502) (DDIM) by Jiaming Song, Chenlin Meng and Stefano Ermon.
The abstract of the paper is the following:
Denoising diffusion probabilistic models (DDPMs) have achieved high quality image generation without adversarial training, yet they require simulating a Markov chain for many steps to produce a sample. To accelerate sampling, we present denoising diffusion implicit models (DDIMs), a more efficient class of iterative implicit probabilistic models with the same training procedure as DDPMs. In DDPMs, the generative process is defined as the reverse of a Markovian diffusion process. We construct a class of non-Markovian diffusion processes that lead to the same training objective, but whose reverse process can be much faster to sample from. We empirically demonstrate that DDIMs can produce high quality samples 10× to 50× faster in terms of wall-clock time compared to DDPMs, allow us to trade off computation for sample quality, and can perform semantically meaningful image interpolation directly in the latent space.
The original codebase of this paper can be found here: [ermongroup/ddim](https://github.com/ermongroup/ddim).
For questions, feel free to contact the author on [tsong.me](https://tsong.me/).
## DDIMScheduler
[[autodoc]] DDIMScheduler

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# Denoising diffusion probabilistic models (DDPM)
## Overview
[Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2006.11239)
(DDPM) by Jonathan Ho, Ajay Jain and Pieter Abbeel proposes the diffusion based model of the same name, but in the context of the 🤗 Diffusers library, DDPM refers to the discrete denoising scheduler from the paper as well as the pipeline.
The abstract of the paper is the following:
We present high quality image synthesis results using diffusion probabilistic models, a class of latent variable models inspired by considerations from nonequilibrium thermodynamics. Our best results are obtained by training on a weighted variational bound designed according to a novel connection between diffusion probabilistic models and denoising score matching with Langevin dynamics, and our models naturally admit a progressive lossy decompression scheme that can be interpreted as a generalization of autoregressive decoding. On the unconditional CIFAR10 dataset, we obtain an Inception score of 9.46 and a state-of-the-art FID score of 3.17. On 256x256 LSUN, we obtain sample quality similar to ProgressiveGAN.
The original paper can be found [here](https://arxiv.org/abs/2010.02502).
## DDPMScheduler
[[autodoc]] DDPMScheduler

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# DPM Discrete Scheduler inspired by Karras et. al paper
## Overview
Inspired by [Karras et. al](https://arxiv.org/abs/2206.00364). Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
## KDPM2DiscreteScheduler
[[autodoc]] KDPM2DiscreteScheduler

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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# DPM Discrete Scheduler with ancestral sampling inspired by Karras et. al paper
## Overview
Inspired by [Karras et. al](https://arxiv.org/abs/2206.00364). Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
## KDPM2AncestralDiscreteScheduler
[[autodoc]] KDPM2AncestralDiscreteScheduler

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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Euler scheduler
## Overview
Euler scheduler (Algorithm 2) from the paper [Elucidating the Design Space of Diffusion-Based Generative Models](https://arxiv.org/abs/2206.00364) by Karras et al. (2022). Based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L51) implementation by Katherine Crowson.
Fast scheduler which often times generates good outputs with 20-30 steps.
## EulerDiscreteScheduler
[[autodoc]] EulerDiscreteScheduler

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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Euler Ancestral scheduler
## Overview
Ancestral sampling with Euler method steps. Based on the original (k-diffusion)[https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72] implementation by Katherine Crowson.
Fast scheduler which often times generates good outputs with 20-30 steps.
## EulerAncestralDiscreteScheduler
[[autodoc]] EulerAncestralDiscreteScheduler

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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Heun scheduler inspired by Karras et. al paper
## Overview
Algorithm 1 of [Karras et. al](https://arxiv.org/abs/2206.00364).
Scheduler ported from @crowsonkb's https://github.com/crowsonkb/k-diffusion library:
All credit for making this scheduler work goes to [Katherine Crowson](https://github.com/crowsonkb/)
## HeunDiscreteScheduler
[[autodoc]] HeunDiscreteScheduler

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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# improved pseudo numerical methods for diffusion models (iPNDM)
## Overview
Original implementation can be found [here](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296).
## IPNDMScheduler
[[autodoc]] IPNDMScheduler

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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Linear multistep scheduler for discrete beta schedules
## Overview
Original implementation can be found [here](https://arxiv.org/abs/2206.00364).
## LMSDiscreteScheduler
[[autodoc]] LMSDiscreteScheduler

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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Multistep DPM-Solver
## Overview
Original paper can be found [here](https://arxiv.org/abs/2206.00927) and the [improved version](https://arxiv.org/abs/2211.01095). The original implementation can be found [here](https://github.com/LuChengTHU/dpm-solver).
## DPMSolverMultistepScheduler
[[autodoc]] DPMSolverMultistepScheduler

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# Schedulers
Diffusers contains multiple pre-built schedule functions for the diffusion process.
## What is a scheduler?
The schedule functions, denoted *Schedulers* in the library take in the output of a trained model, a sample which the diffusion process is iterating on, and a timestep to return a denoised sample. That's why schedulers may also be called *Samplers* in other diffusion models implementations.
- Schedulers define the methodology for iteratively adding noise to an image or for updating a sample based on model outputs.
- adding noise in different manners represent the algorithmic processes to train a diffusion model by adding noise to images.
- for inference, the scheduler defines how to update a sample based on an output from a pretrained model.
- Schedulers are often defined by a *noise schedule* and an *update rule* to solve the differential equation solution.
### Discrete versus continuous schedulers
All schedulers take in a timestep to predict the updated version of the sample being diffused.
The timesteps dictate where in the diffusion process the step is, where data is generated by iterating forward in time and inference is executed by propagating backwards through timesteps.
Different algorithms use timesteps that can be discrete (accepting `int` inputs), such as the [`DDPMScheduler`] or [`PNDMScheduler`], or continuous (accepting `float` inputs), such as the score-based schedulers [`ScoreSdeVeScheduler`] or [`ScoreSdeVpScheduler`].
## Designing Re-usable schedulers
The core design principle between the schedule functions is to be model, system, and framework independent.
This allows for rapid experimentation and cleaner abstractions in the code, where the model prediction is separated from the sample update.
To this end, the design of schedulers is such that:
- Schedulers can be used interchangeably between diffusion models in inference to find the preferred trade-off between speed and generation quality.
- Schedulers are currently by default in PyTorch, but are designed to be framework independent (partial Jax support currently exists).
## Schedulers Summary
The following table summarizes all officially supported schedulers, their corresponding paper
| Scheduler | Paper |
|---|---|
| [ddim](./ddim) | [**Denoising Diffusion Implicit Models**](https://arxiv.org/abs/2010.02502) |
| [ddpm](./ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) |
| [singlestep_dpm_solver](./singlestep_dpm_solver) | [**Singlestep DPM-Solver**](https://arxiv.org/abs/2206.00927) |
| [multistep_dpm_solver](./multistep_dpm_solver) | [**Multistep DPM-Solver**](https://arxiv.org/abs/2206.00927) |
| [heun](./heun) | [**Heun scheduler inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
| [dpm_discrete](./dpm_discrete) | [**DPM Discrete Scheduler inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
| [dpm_discrete_ancestral](./dpm_discrete_ancestral) | [**DPM Discrete Scheduler with ancestral sampling inspired by Karras et. al paper**](https://arxiv.org/abs/2206.00364) |
| [stochastic_karras_ve](./stochastic_karras_ve) | [**Variance exploding, stochastic sampling from Karras et. al**](https://arxiv.org/abs/2206.00364) |
| [lms_discrete](./lms_discrete) | [**Linear multistep scheduler for discrete beta schedules**](https://arxiv.org/abs/2206.00364) |
| [pndm](./pndm) | [**Pseudo numerical methods for diffusion models (PNDM)**](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181) |
| [score_sde_ve](./score_sde_ve) | [**variance exploding stochastic differential equation (VE-SDE) scheduler**](https://arxiv.org/abs/2011.13456) |
| [ipndm](./ipndm) | [**improved pseudo numerical methods for diffusion models (iPNDM)**](https://github.com/crowsonkb/v-diffusion-pytorch/blob/987f8985e38208345c1959b0ea767a625831cc9b/diffusion/sampling.py#L296) |
| [score_sde_vp](./score_sde_vp) | [**Variance preserving stochastic differential equation (VP-SDE) scheduler**](https://arxiv.org/abs/2011.13456) |
| [euler](./euler) | [**Euler scheduler**](https://arxiv.org/abs/2206.00364) |
| [euler_ancestral](./euler_ancestral) | [**Euler Ancestral scheduler**](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L72) |
| [vq_diffusion](./vq_diffusion) | [**VQDiffusionScheduler**](https://arxiv.org/abs/2111.14822) |
| [repaint](./repaint) | [**RePaint scheduler**](https://arxiv.org/abs/2201.09865) |
## API
The core API for any new scheduler must follow a limited structure.
- Schedulers should provide one or more `def step(...)` functions that should be called to update the generated sample iteratively.
- Schedulers should provide a `set_timesteps(...)` method that configures the parameters of a schedule function for a specific inference task.
- Schedulers should be framework-specific.
The base class [`SchedulerMixin`] implements low level utilities used by multiple schedulers.
### SchedulerMixin
[[autodoc]] SchedulerMixin
### SchedulerOutput
The class [`SchedulerOutput`] contains the outputs from any schedulers `step(...)` call.
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput

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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Pseudo numerical methods for diffusion models (PNDM)
## Overview
Original implementation can be found [here](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L181).
## PNDMScheduler
[[autodoc]] PNDMScheduler

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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
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# RePaint scheduler
## Overview
DDPM-based inpainting scheduler for unsupervised inpainting with extreme masks.
Intended for use with [`RePaintPipeline`].
Based on the paper [RePaint: Inpainting using Denoising Diffusion Probabilistic Models](https://arxiv.org/abs/2201.09865)
and the original implementation by Andreas Lugmayr et al.: https://github.com/andreas128/RePaint
## RePaintScheduler
[[autodoc]] RePaintScheduler

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# variance exploding stochastic differential equation (VE-SDE) scheduler
## Overview
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
## ScoreSdeVeScheduler
[[autodoc]] ScoreSdeVeScheduler

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Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Variance preserving stochastic differential equation (VP-SDE) scheduler
## Overview
Original paper can be found [here](https://arxiv.org/abs/2011.13456).
<Tip warning={true}>
Score SDE-VP is under construction.
</Tip>
## ScoreSdeVpScheduler
[[autodoc]] schedulers.scheduling_sde_vp.ScoreSdeVpScheduler

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<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
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Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
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# Singlestep DPM-Solver
## Overview
Original paper can be found [here](https://arxiv.org/abs/2206.00927) and the [improved version](https://arxiv.org/abs/2211.01095). The original implementation can be found [here](https://github.com/LuChengTHU/dpm-solver).
## DPMSolverSinglestepScheduler
[[autodoc]] DPMSolverSinglestepScheduler

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# Variance exploding, stochastic sampling from Karras et. al
## Overview
Original paper can be found [here](https://arxiv.org/abs/2206.00364).
## KarrasVeScheduler
[[autodoc]] KarrasVeScheduler

View File

@@ -0,0 +1,20 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# VQDiffusionScheduler
## Overview
Original paper can be found [here](https://arxiv.org/abs/2111.14822)
## VQDiffusionScheduler
[[autodoc]] VQDiffusionScheduler

View File

@@ -18,12 +18,12 @@ specific language governing permissions and limitations under the License.
# 🧨 Diffusers
🤗 Diffusers provides pretrained vision diffusion models, and serves as a modular toolbox for inference and training.
🤗 Diffusers provides pretrained vision and audio diffusion models, and serves as a modular toolbox for inference and training.
More precisely, 🤗 Diffusers offers:
- State-of-the-art diffusion pipelines that can be run in inference with just a couple of lines of code (see [**Using Diffusers**](./using-diffusers/conditional_image_generation)) or have a look at [**Pipelines**](#pipelines) to get an overview of all supported pipelines and their corresponding papers.
- Various noise schedulers that can be used interchangeably for the preferred speed vs. quality trade-off in inference. For more information see [**Schedulers**](./api/schedulers).
- Various noise schedulers that can be used interchangeably for the preferred speed vs. quality trade-off in inference. For more information see [**Schedulers**](./api/schedulers/overview).
- Multiple types of models, such as UNet, can be used as building blocks in an end-to-end diffusion system. See [**Models**](./api/models) for more details
- Training examples to show how to train the most popular diffusion model tasks. For more information see [**Training**](./training/overview).
@@ -35,6 +35,7 @@ available a colab notebook to directly try them out.
| Pipeline | Paper | Tasks | Colab
|---|---|:---:|:---:|
| [alt_diffusion](./api/pipelines/alt_diffusion) | [**AltDiffusion**](https://arxiv.org/abs/2211.06679) | Image-to-Image Text-Guided Generation |
| [audio_diffusion](./api/pipelines/audio_diffusion) | [**Audio Diffusion**](https://github.com/teticio/audio-diffusion.git) | Unconditional Audio Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/teticio/audio-diffusion/blob/master/notebooks/audio_diffusion_pipeline.ipynb)
| [cycle_diffusion](./api/pipelines/cycle_diffusion) | [**Cycle Diffusion**](https://arxiv.org/abs/2210.05559) | Image-to-Image Text-Guided Generation |
| [dance_diffusion](./api/pipelines/dance_diffusion) | [**Dance Diffusion**](https://github.com/williamberman/diffusers.git) | Unconditional Audio Generation |
| [ddpm](./api/pipelines/ddpm) | [**Denoising Diffusion Probabilistic Models**](https://arxiv.org/abs/2006.11239) | Unconditional Image Generation |
@@ -42,6 +43,7 @@ available a colab notebook to directly try them out.
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Text-to-Image Generation |
| [latent_diffusion](./api/pipelines/latent_diffusion) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)| Super Resolution Image-to-Image |
| [latent_diffusion_uncond](./api/pipelines/latent_diffusion_uncond) | [**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752) | Unconditional Image Generation |
| [paint_by_example](./api/pipelines/paint_by_example) | [**Paint by Example: Exemplar-based Image Editing with Diffusion Models**](https://arxiv.org/abs/2211.13227) | Image-Guided Image Inpainting |
| [pndm](./api/pipelines/pndm) | [**Pseudo Numerical Methods for Diffusion Models on Manifolds**](https://arxiv.org/abs/2202.09778) | Unconditional Image Generation |
| [score_sde_ve](./api/pipelines/score_sde_ve) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
| [score_sde_vp](./api/pipelines/score_sde_vp) | [**Score-Based Generative Modeling through Stochastic Differential Equations**](https://openreview.net/forum?id=PxTIG12RRHS) | Unconditional Image Generation |
@@ -53,6 +55,7 @@ available a colab notebook to directly try them out.
| [stable_diffusion_2](./api/pipelines/stable_diffusion_2) | [**Stable Diffusion 2**](https://stability.ai/blog/stable-diffusion-v2-release) | Text-Guided Super Resolution Image-to-Image |
| [stable_diffusion_safe](./api/pipelines/stable_diffusion_safe) | [**Safe Stable Diffusion**](https://arxiv.org/abs/2211.05105) | Text-Guided Generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/ml-research/safe-latent-diffusion/blob/main/examples/Safe%20Latent%20Diffusion.ipynb)
| [stochastic_karras_ve](./api/pipelines/stochastic_karras_ve) | [**Elucidating the Design Space of Diffusion-Based Generative Models**](https://arxiv.org/abs/2206.00364) | Unconditional Image Generation |
| [unclip](./api/pipelines/unclip) | [Hierarchical Text-Conditional Image Generation with CLIP Latents](https://arxiv.org/abs/2204.06125) | Text-to-Image Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Text-to-Image Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Image Variations Generation |
| [versatile_diffusion](./api/pipelines/versatile_diffusion) | [Versatile Diffusion: Text, Images and Variations All in One Diffusion Model](https://arxiv.org/abs/2211.08332) | Dual Image and Text Guided Generation |

View File

@@ -120,3 +120,25 @@ git pull
```
Your Python environment will find the `main` version of 🤗 Diffusers on the next run.
## Notice on telemetry logging
Our library gathers telemetry information during `from_pretrained()` requests.
This data includes the version of Diffusers and PyTorch/Flax, the requested model or pipeline class,
and the path to a pretrained checkpoint if it is hosted on the Hub.
This usage data helps us debug issues and prioritize new features.
Telemetry is only sent when loading models and pipelines from the HuggingFace Hub,
and is not collected during local usage.
We understand that not everyone wants to share additional information, and we respect your privacy,
so you can disable telemetry collection by setting the `DISABLE_TELEMETRY` environment variable from your terminal:
On Linux/MacOS:
```bash
export DISABLE_TELEMETRY=YES
```
On Windows:
```bash
set DISABLE_TELEMETRY=YES
```

View File

@@ -12,7 +12,9 @@ specific language governing permissions and limitations under the License.
# Memory and speed
We present some techniques and ideas to optimize 🤗 Diffusers _inference_ for memory or speed.
We present some techniques and ideas to optimize 🤗 Diffusers _inference_ for memory or speed. As a general rule, we recommend the use of [xFormers](https://github.com/facebookresearch/xformers) for memory efficient attention, please see the recommended [installation instructions](xformers).
We'll discuss how the following settings impact performance and memory.
| | Latency | Speedup |
| ---------------- | ------- | ------- |
@@ -77,7 +79,7 @@ To save more GPU memory and get even more speed, you can load and run the model
```Python
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
@@ -105,7 +107,7 @@ from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
@@ -132,7 +134,7 @@ from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
@@ -157,7 +159,7 @@ from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
@@ -177,7 +179,7 @@ from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")
@@ -232,7 +234,6 @@ def generate_inputs():
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
).to("cuda")
unet = pipe.unet
@@ -296,7 +297,6 @@ class UNet2DConditionOutput:
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
).to("cuda")
@@ -322,7 +322,9 @@ with torch.inference_mode():
## Memory Efficient Attention
Recent work on optimizing the bandwitdh in the attention block have generated huge speed ups and gains in GPU memory usage. The most recent being Flash Attention (from @tridao, [code](https://github.com/HazyResearch/flash-attention), [paper](https://arxiv.org/pdf/2205.14135.pdf)) .
Recent work on optimizing the bandwitdh in the attention block has generated huge speed ups and gains in GPU memory usage. The most recent being Flash Attention from @tridao: [code](https://github.com/HazyResearch/flash-attention), [paper](https://arxiv.org/pdf/2205.14135.pdf).
Here are the speedups we obtain on a few Nvidia GPUs when running the inference at 512x512 with a batch size of 1 (one prompt):
| GPU | Base Attention FP16 | Memory Efficient Attention FP16 |
@@ -338,14 +340,13 @@ Here are the speedups we obtain on a few Nvidia GPUs when running the inference
To leverage it just make sure you have:
- PyTorch > 1.12
- Cuda available
- Installed the [xformers](https://github.com/facebookresearch/xformers) library
- [Installed the xformers library](xformers).
```python
from diffusers import StableDiffusionPipeline
import torch
pipe = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
revision="fp16",
torch_dtype=torch.float16,
).to("cuda")

View File

@@ -0,0 +1,70 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# How to use Stable Diffusion on Habana Gaudi
🤗 Diffusers is compatible with Habana Gaudi through 🤗 [Optimum Habana](https://huggingface.co/docs/optimum/habana/usage_guides/stable_diffusion).
## Requirements
- Optimum Habana 1.3 or later, [here](https://huggingface.co/docs/optimum/habana/installation) is how to install it.
- SynapseAI 1.7.
## Inference Pipeline
To generate images with Stable Diffusion 1 and 2 on Gaudi, you need to instantiate two instances:
- A pipeline with [`GaudiStableDiffusionPipeline`](https://huggingface.co/docs/optimum/habana/package_reference/stable_diffusion_pipeline). This pipeline supports *text-to-image generation*.
- A scheduler with [`GaudiDDIMScheduler`](https://huggingface.co/docs/optimum/habana/package_reference/stable_diffusion_pipeline#optimum.habana.diffusers.GaudiDDIMScheduler). This scheduler has been optimized for Habana Gaudi.
When initializing the pipeline, you have to specify `use_habana=True` to deploy it on HPUs.
Furthermore, in order to get the fastest possible generations you should enable **HPU graphs** with `use_hpu_graphs=True`.
Finally, you will need to specify a [Gaudi configuration](https://huggingface.co/docs/optimum/habana/package_reference/gaudi_config) which can be downloaded from the [Hugging Face Hub](https://huggingface.co/Habana).
```python
from optimum.habana import GaudiConfig
from optimum.habana.diffusers import GaudiDDIMScheduler, GaudiStableDiffusionPipeline
model_name = "stabilityai/stable-diffusion-2-base"
scheduler = GaudiDDIMScheduler.from_pretrained(model_name, subfolder="scheduler")
pipeline = GaudiStableDiffusionPipeline.from_pretrained(
model_name,
scheduler=scheduler,
use_habana=True,
use_hpu_graphs=True,
gaudi_config="Habana/stable-diffusion",
)
```
You can then call the pipeline to generate images by batches from one or several prompts:
```python
outputs = pipeline(
prompt=[
"High quality photo of an astronaut riding a horse in space",
"Face of a yellow cat, high resolution, sitting on a park bench",
],
num_images_per_prompt=10,
batch_size=4,
)
```
For more information, check out Optimum Habana's [documentation](https://huggingface.co/docs/optimum/habana/usage_guides/stable_diffusion) and the [example](https://github.com/huggingface/optimum-habana/tree/main/examples/stable-diffusion) provided in the official Github repository.
## Benchmark
Here are the latencies for Habana Gaudi 1 and Gaudi 2 with the [Habana/stable-diffusion](https://huggingface.co/Habana/stable-diffusion) Gaudi configuration (mixed precision bf16/fp32):
| | Latency | Batch size |
| ------- |:-------:|:----------:|
| Gaudi 1 | 4.37s | 4/8 |
| Gaudi 2 | 1.19s | 4/8 |

View File

@@ -0,0 +1,26 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Installing xFormers
We recommend the use of [xFormers](https://github.com/facebookresearch/xformers) for both inference and training. In our tests, the optimizations performed in the attention blocks allow for both faster speed and reduced memory consumption.
Installing xFormers has historically been a bit involved, as binary distributions were not always up to date. Fortunately, the project has [very recently](https://github.com/facebookresearch/xformers/pull/591) integrated a process to build pip wheels as part of the project's continuous integration, so this should improve a lot starting from xFormers version 0.0.16.
Until xFormers 0.0.16 is deployed, you can install pip wheels using [`TestPyPI`](https://test.pypi.org/project/formers/). These are the steps that worked for us in a Linux computer to install xFormers version 0.0.15:
```bash
pip install pyre-extensions==0.0.23
pip install -i https://test.pypi.org/simple/ formers==0.0.15.dev376
```
We'll update these instructions when the wheels are published to the official PyPI repository.

View File

@@ -18,9 +18,12 @@ Whether you're a developer or an everyday user, this quick tour will help you ge
Before you begin, make sure you have all the necessary libraries installed:
```bash
pip install --upgrade diffusers
pip install --upgrade diffusers accelerate transformers
```
- [`accelerate`](https://huggingface.co/docs/accelerate/index) speeds up model loading for inference and training
- [`transformers`](https://huggingface.co/docs/transformers/index) is required to run the most popular diffusion models, such as [Stable Diffusion](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion)
## DiffusionPipeline
The [`DiffusionPipeline`] is the easiest way to use a pre-trained diffusion system for inference. You can use the [`DiffusionPipeline`] out-of-the-box for many tasks across different modalities. Take a look at the table below for some supported tasks:
@@ -29,19 +32,26 @@ The [`DiffusionPipeline`] is the easiest way to use a pre-trained diffusion syst
|------------------------------|--------------------------------------------------------------------------------------------------------------|-----------------|
| Unconditional Image Generation | generate an image from gaussian noise | [unconditional_image_generation](./using-diffusers/unconditional_image_generation`) |
| Text-Guided Image Generation | generate an image given a text prompt | [conditional_image_generation](./using-diffusers/conditional_image_generation) |
| Text-Guided Image-to-Image Translation | generate an image given an original image and a text prompt | [img2img](./using-diffusers/img2img) |
| Text-Guided Image-to-Image Translation | adapt an image guided by a text prompt | [img2img](./using-diffusers/img2img) |
| Text-Guided Image-Inpainting | fill the masked part of an image given the image, the mask and a text prompt | [inpaint](./using-diffusers/inpaint) |
| Text-Guided Depth-to-Image Translation | adapt parts of an image guided by a text prompt while preserving structure via depth estimation | [depth2image](./using-diffusers/depth2image) |
For more in-detail information on how diffusion pipelines function for the different tasks, please have a look at the [**Using Diffusers**](./using-diffusers/overview) section.
As an example, start by creating an instance of [`DiffusionPipeline`] and specify which pipeline checkpoint you would like to download.
You can use the [`DiffusionPipeline`] for any [Diffusers' checkpoint](https://huggingface.co/models?library=diffusers&sort=downloads).
In this guide though, you'll use [`DiffusionPipeline`] for text-to-image generation with [Latent Diffusion](https://huggingface.co/CompVis/ldm-text2im-large-256):
In this guide though, you'll use [`DiffusionPipeline`] for text-to-image generation with [Stable Diffusion](https://huggingface.co/CompVis/stable-diffusion).
For [Stable Diffusion](https://huggingface.co/CompVis/stable-diffusion), please carefully read its [license](https://huggingface.co/spaces/CompVis/stable-diffusion-license) before running the model.
This is due to the improved image generation capabilities of the model and the potentially harmful content that could be produced with it.
Please, head over to your stable diffusion model of choice, *e.g.* [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5), and read the license.
You can load the model as follows:
```python
>>> from diffusers import DiffusionPipeline
>>> pipeline = DiffusionPipeline.from_pretrained("CompVis/ldm-text2im-large-256")
>>> pipeline = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
```
The [`DiffusionPipeline`] downloads and caches all modeling, tokenization, and scheduling components.
@@ -66,40 +76,14 @@ You can save the image by simply calling:
>>> image.save("image_of_squirrel_painting.png")
```
More advanced models, like [Stable Diffusion](https://huggingface.co/CompVis/stable-diffusion) require you to accept a [license](https://huggingface.co/spaces/CompVis/stable-diffusion-license) before running the model.
This is due to the improved image generation capabilities of the model and the potentially harmful content that could be produced with it.
Please, head over to your stable diffusion model of choice, *e.g.* [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5), read the license carefully and tick the checkbox if you agree.
You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section of the documentation](https://huggingface.co/docs/hub/security-tokens).
Having "click-accepted" the license, you can save your token:
```python
AUTH_TOKEN = "<please-fill-with-your-token>"
```
You can then load [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5)
just like we did before only that now you need to pass your `AUTH_TOKEN`:
```python
>>> from diffusers import DiffusionPipeline
>>> pipeline = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", use_auth_token=AUTH_TOKEN)
```
If you do not pass your authentication token you will see that the diffusion system will not be correctly
downloaded. Forcing the user to pass an authentication token ensures that it can be verified that the
user has indeed read and accepted the license, which also means that an internet connection is required.
**Note**: If you do not want to be forced to pass an authentication token, you can also simply download
the weights locally via:
**Note**: You can also use the pipeline locally by downloading the weights via:
```
git lfs install
git clone https://huggingface.co/runwayml/stable-diffusion-v1-5
```
and then load locally saved weights into the pipeline. This way, you do not need to pass an authentication
token. Assuming that `"./stable-diffusion-v1-5"` is the local path to the cloned stable-diffusion-v1-5 repo,
you can also load the pipeline as follows:
and then loading the saved weights into the pipeline.
```python
>>> pipeline = DiffusionPipeline.from_pretrained("./stable-diffusion-v1-5")
@@ -113,7 +97,7 @@ Running the pipeline is then identical to the code above as it's the same model
>>> image.save("image_of_squirrel_painting.png")
```
Diffusion systems can be used with multiple different [schedulers](./api/schedulers) each with their
Diffusion systems can be used with multiple different [schedulers](./api/schedulers/overview) each with their
pros and cons. By default, Stable Diffusion runs with [`PNDMScheduler`], but it's very simple to
use a different scheduler. *E.g.* if you would instead like to use the [`EulerDiscreteScheduler`] scheduler,
you could use it as follows:
@@ -121,7 +105,7 @@ you could use it as follows:
```python
>>> from diffusers import EulerDiscreteScheduler
>>> pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", use_auth_token=AUTH_TOKEN)
>>> pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
>>> # change scheduler to Euler
>>> pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)

View File

@@ -21,8 +21,6 @@ The [Dreambooth training script](https://github.com/huggingface/diffusers/tree/m
<Tip warning={true}>
<!-- TODO: replace with our blog when it's done -->
Dreambooth fine-tuning is very sensitive to hyperparameters and easy to overfit. We recommend you take a look at our [in-depth analysis](https://huggingface.co/blog/dreambooth) with recommended settings for different subjects, and go from there.
</Tip>
@@ -38,23 +36,17 @@ pip install git+https://github.com/huggingface/diffusers
pip install -U -r diffusers/examples/dreambooth/requirements.txt
```
Then initialize and configure a [🤗 Accelerate](https://github.com/huggingface/accelerate/) environment with:
xFormers is not part of the training requirements, but [we recommend you install it if you can](../optimization/xformers). It could make your training faster and less memory intensive.
After all dependencies have been set up you can configure a [🤗 Accelerate](https://github.com/huggingface/accelerate/) environment with:
```bash
accelerate config
```
You need to accept the model license before downloading or using the weights. In this example we'll use model version `v1-4`, so you'll need to visit [its card](https://huggingface.co/CompVis/stable-diffusion-v1-4), read the license and tick the checkbox if you agree.
In this example we'll use model version `v1-4`, so please visit [its card](https://huggingface.co/CompVis/stable-diffusion-v1-4) and carefully read the license before proceeding.
You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section of the documentation](https://huggingface.co/docs/hub/security-tokens).
Run the following command to authenticate your token
```bash
huggingface-cli login
```
If you have already cloned the repo, then you won't need to go through these steps. Instead, you can pass the path to your local checkout to the training script and it will be loaded from there.
The command below will download and cache the model weights from the Hub because we use the model's Hub id `CompVis/stable-diffusion-v1-4`. You may also clone the repo locally and use the local path in your system where the checkout was saved.
### Dog toy example
@@ -111,6 +103,59 @@ accelerate launch train_dreambooth.py \
--max_train_steps=800
```
### Saving checkpoints while training
It's easy to overfit while training with Dreambooth, so sometimes it's useful to save regular checkpoints during the process. One of the intermediate checkpoints might work better than the final model! To use this feature you need to pass the following argument to the training script:
```bash
--checkpointing_steps=500
```
This will save the full training state in subfolders of your `output_dir`. Subfolder names begin with the prefix `checkpoint-`, and then the number of steps performed so far; for example: `checkpoint-1500` would be a checkpoint saved after 1500 training steps.
#### Resuming training from a saved checkpoint
If you want to resume training from any of the saved checkpoints, you can pass the argument `--resume_from_checkpoint` and then indicate the name of the checkpoint you want to use. You can also use the special string `"latest"` to resume from the last checkpoint saved (i.e., the one with the largest number of steps). For example, the following would resume training from the checkpoint saved after 1500 steps:
```bash
--resume_from_checkpoint="checkpoint-1500"
```
This would be a good opportunity to tweak some of your hyperparameters if you wish.
#### Performing inference using a saved checkpoint
Saved checkpoints are stored in a format suitable for resuming training. They not only include the model weights, but also the state of the optimizer, data loaders and learning rate.
You can use a checkpoint for inference, but first you need to convert it to an inference pipeline. This is how you could do it:
```python
from accelerate import Accelerator
from diffusers import DiffusionPipeline
# Load the pipeline with the same arguments (model, revision) that were used for training
model_id = "CompVis/stable-diffusion-v1-4"
pipeline = DiffusionPipeline.from_pretrained(model_id)
accelerator = Accelerator()
# Use text_encoder if `--train_text_encoder` was used for the initial training
unet, text_encoder = accelerator.prepare(pipeline.unet, pipeline.text_encoder)
# Restore state from a checkpoint path. You have to use the absolute path here.
accelerator.load_state("/sddata/dreambooth/daruma-v2-1/checkpoint-100")
# Rebuild the pipeline with the unwrapped models (assignment to .unet and .text_encoder should work too)
pipeline = DiffusionPipeline.from_pretrained(
model_id,
unet=accelerator.unwrap_model(unet),
text_encoder=accelerator.unwrap_model(text_encoder),
)
# Perform inference, or save, or push to the hub
pipeline.save_pretrained("dreambooth-pipeline")
```
### Training on a 16GB GPU
With the help of gradient checkpointing and the 8-bit optimizer from [bitsandbytes](https://github.com/TimDettmers/bitsandbytes), it's possible to train dreambooth on a 16GB GPU.

View File

@@ -38,6 +38,7 @@ Training examples show how to pretrain or fine-tune diffusion models for a varie
- [Text Inversion](./text_inversion)
- [Dreambooth](./dreambooth)
If possible, please [install xFormers](../optimization/xformers) for memory efficient attention. This could help make your training faster and less memory intensive.
| Task | 🤗 Accelerate | 🤗 Datasets | Colab
|---|---|:---:|:---:|

View File

@@ -12,5 +12,5 @@ specific language governing permissions and limitations under the License.
# Using Diffusers for audio
The [`DanceDiffusionPipeline`] can be used to generate audio rapidly!
More coming soon!
[`DanceDiffusionPipeline`] and [`AudioDiffusionPipeline`] can be used to generate
audio rapidly! More coming soon!

View File

@@ -58,7 +58,6 @@ guided_pipeline = DiffusionPipeline.from_pretrained(
custom_pipeline="clip_guided_stable_diffusion",
clip_model=clip_model,
feature_extractor=feature_extractor,
revision="fp16",
torch_dtype=torch.float16,
)
guided_pipeline.enable_attention_slicing()
@@ -113,7 +112,6 @@ import torch
pipe = DiffusionPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
revision="fp16",
torch_dtype=torch.float16,
safety_checker=None, # Very important for videos...lots of false positives while interpolating
custom_pipeline="interpolate_stable_diffusion",
@@ -159,7 +157,6 @@ pipe = DiffusionPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
custom_pipeline="stable_diffusion_mega",
torch_dtype=torch.float16,
revision="fp16",
)
pipe.to("cuda")
pipe.enable_attention_slicing()
@@ -177,7 +174,7 @@ init_image = download_image(
prompt = "A fantasy landscape, trending on artstation"
images = pipe.img2img(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
images = pipe.img2img(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
### Inpainting
@@ -187,7 +184,7 @@ init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
prompt = "a cat sitting on a bench"
images = pipe.inpaint(prompt=prompt, init_image=init_image, mask_image=mask_image, strength=0.75).images
images = pipe.inpaint(prompt=prompt, image=init_image, mask_image=mask_image, strength=0.75).images
```
As shown above this one pipeline can run all both "text-to-image", "image-to-image", and "inpainting" in one pipeline.
@@ -204,7 +201,7 @@ from diffusers import DiffusionPipeline
import torch
pipe = DiffusionPipeline.from_pretrained(
"hakurei/waifu-diffusion", custom_pipeline="lpw_stable_diffusion", revision="fp16", torch_dtype=torch.float16
"hakurei/waifu-diffusion", custom_pipeline="lpw_stable_diffusion", torch_dtype=torch.float16
)
pipe = pipe.to("cuda")
@@ -268,7 +265,7 @@ diffuser_pipeline = DiffusionPipeline.from_pretrained(
custom_pipeline="speech_to_image_diffusion",
speech_model=model,
speech_processor=processor,
revision="fp16",
torch_dtype=torch.float16,
)

View File

@@ -0,0 +1,35 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Text-Guided Image-to-Image Generation
The [`StableDiffusionDepth2ImgPipeline`] lets you pass a text prompt and an initial image to condition the generation of new images as well as a `depth_map` to preserve the images' structure. If no `depth_map` is provided, the pipeline will automatically predict the depth via an integrated depth-estimation model.
```python
import torch
import requests
from PIL import Image
from diffusers import StableDiffusionDepth2ImgPipeline
pipe = StableDiffusionDepth2ImgPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-depth",
torch_dtype=torch.float16,
).to("cuda")
url = "http://images.cocodataset.org/val2017/000000039769.jpg"
init_image = Image.open(requests.get(url, stream=True).raw)
prompt = "two tigers"
n_prompt = "bad, deformed, ugly, bad anatomy"
image = pipe(prompt=prompt, image=init_image, negative_prompt=n_prompt, strength=0.7).images[0]
```

View File

@@ -24,9 +24,9 @@ from diffusers import StableDiffusionImg2ImgPipeline
# load the pipeline
device = "cuda"
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", revision="fp16", torch_dtype=torch.float16
).to(device)
pipe = StableDiffusionImg2ImgPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16).to(
device
)
# let's download an initial image
url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
@@ -37,7 +37,7 @@ init_image.thumbnail((768, 768))
prompt = "A fantasy landscape, trending on artstation"
images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
images[0].save("fantasy_landscape.png")
```

View File

@@ -42,7 +42,6 @@ mask_image = download_image(mask_url).resize((512, 512))
pipe = StableDiffusionInpaintPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
revision="fp16",
torch_dtype=torch.float16,
)
pipe = pipe.to("cuda")

View File

@@ -14,7 +14,8 @@ specific language governing permissions and limitations under the License.
Diffusers is in the process of expanding to modalities other than images.
Currently, one example is for [molecule conformation](https://www.nature.com/subjects/molecular-conformation#:~:text=Definition,to%20changes%20in%20their%20environment.) generation.
* Generate conformations in Colab [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/geodiff_molecule_conformation.ipynb)
Example type | Colab | Pipeline |
:-------------------------:|:-------------------------:|:-------------------------:|
[Molecule conformation](https://www.nature.com/subjects/molecular-conformation#:~:text=Definition,to%20changes%20in%20their%20environment.) generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/geodiff_molecule_conformation.ipynb) | ❌
More coming soon!

View File

@@ -0,0 +1,73 @@
<!--Copyright 2022 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Re-using seeds for fast prompt engineering
A common use case when generating images is to generate a batch of images, select one image and improve it with a better, more detailed prompt in a second run.
To do this, one needs to make each generated image of the batch deterministic.
Images are generated by denoising gaussian random noise which can be instantiated by passing a [torch generator](https://pytorch.org/docs/stable/generated/torch.Generator.html#generator).
Now, for batched generation, we need to make sure that every single generated image in the batch is tied exactly to one seed. In 🧨 Diffusers, this can be achieved by not passing one `generator`, but a list
of `generators` to the pipeline.
Let's go through an example using [`runwayml/stable-diffusion-v1-5`](runwayml/stable-diffusion-v1-5).
We want to generate several versions of the prompt:
```py
prompt = "Labrador in the style of Vermeer"
```
Let's load the pipeline
```python
>>> from diffusers import DiffusionPipeline
>>> pipe = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
>>> pipe = pipe.to("cuda")
```
Now, let's define 4 different generators, since we would like to reproduce a certain image. We'll use seeds `0` to `3` to create our generators.
```python
>>> import torch
>>> generator = [torch.Generator(device="cuda").manual_seed(i) for i in range(4)]
```
Let's generate 4 images:
```python
>>> images = pipe(prompt, generator=generator, num_images_per_prompt=4).images
>>> images
```
![img](https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/reusabe_seeds.jpg)
Ok, the last images has some double eyes, but the first image looks good!
Let's try to make the prompt a bit better **while keeping the first seed**
so that the images are similar to the first image.
```python
prompt = [prompt + t for t in [", highly realistic", ", artsy", ", trending", ", colorful"]]
generator = [torch.Generator(device="cuda").manual_seed(0) for i in range(4)]
```
We create 4 generators with seed `0`, which is the first seed we used before.
Let's run the pipeline again.
```python
>>> images = pipe(prompt, generator=generator).images
>>> images
```
![img](https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/reusabe_seeds_2.jpg)

View File

@@ -13,6 +13,13 @@ specific language governing permissions and limitations under the License.
# Using Diffusers for reinforcement learning
Support for one RL model and related pipelines is included in the `experimental` source of diffusers.
More models and examples coming soon!
To try some of this in colab, please look at the following example:
* Model-based reinforcement learning on Colab [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/reinforcement_learning_with_diffusers.ipynb) ![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)
# Diffuser Value-guided Planning
You can run the model from [*Planning with Diffusion for Flexible Behavior Synthesis*](https://arxiv.org/abs/2205.09991) with Diffusers.
The script is located in the [RL Examples](https://github.com/huggingface/diffusers/tree/main/examples/rl) folder.
Or, run this example in Colab [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/reinforcement_learning_with_diffusers.ipynb)
[[autodoc]] diffusers.experimental.ValueGuidedRLPipeline

View File

@@ -13,7 +13,7 @@ specific language governing permissions and limitations under the License.
# Schedulers
Diffusion pipelines are inherently a collection of diffusion models and schedulers that are partly independent from each other. This means that one is able to switch out parts of the pipeline to better customize
a pipeline to one's use case. The best example of this are the [Schedulers](../api/schedulers.mdx).
a pipeline to one's use case. The best example of this are the [Schedulers](../api/schedulers/overview.mdx).
Whereas diffusion models usually simply define the forward pass from noise to a less noisy sample,
schedulers define the whole denoising process, *i.e.*:

View File

@@ -52,6 +52,10 @@ For such examples, we are more lenient regarding the philosophy defined above an
Examples that are useful for the community, but are either not yet deemed popular or not yet following our above philosophy should go into the [community examples](https://github.com/huggingface/diffusers/tree/main/examples/community) folder. The community folder therefore includes training examples and inference pipelines.
**Note**: Community examples can be a [great first contribution](https://github.com/huggingface/diffusers/issues?q=is%3Aopen+is%3Aissue+label%3A%22good+first+issue%22) to show to the community how you like to use `diffusers` 🪄.
## Research Projects
We also provide **research_projects** examples that are maintained by the community as defined in the respective research project folders. These examples are useful and offer the extended capabilities which are complementary to the official examples. You may refer to [research_projects](https://github.com/huggingface/diffusers/tree/main/examples/research_projects) for details.
## Important note
To make sure you can successfully run the latest versions of the example scripts, you have to **install the library from source** and install some example-specific requirements. To do this, execute the following steps in a new virtual environment:

View File

@@ -23,6 +23,8 @@ If a community doesn't work as expected, please open an issue and ping the autho
| Text Based Inpainting Stable Diffusion | Stable Diffusion Inpainting Pipeline that enables passing a text prompt to generate the mask for inpainting| [Text Based Inpainting Stable Diffusion](#image-to-image-inpainting-stable-diffusion) | - | [Dhruv Karan](https://github.com/unography) |
| Bit Diffusion | Diffusion on discrete data | [Bit Diffusion](#bit-diffusion) | - |[Stuti R.](https://github.com/kingstut) |
| K-Diffusion Stable Diffusion | Run Stable Diffusion with any of [K-Diffusion's samplers](https://github.com/crowsonkb/k-diffusion/blob/master/k_diffusion/sampling.py) | [Stable Diffusion with K Diffusion](#stable-diffusion-with-k-diffusion) | - | [Patrick von Platen](https://github.com/patrickvonplaten/) |
| Checkpoint Merger Pipeline | Diffusion Pipeline that enables merging of saved model checkpoints | [Checkpoint Merger Pipeline](#checkpoint-merger-pipeline) | - | [Naga Sai Abhinay Devarinti](https://github.com/Abhinay1997/) |
Stable Diffusion v1.1-1.4 Comparison | Run all 4 model checkpoints for Stable Diffusion and compare their results together | [Stable Diffusion Comparison](#stable-diffusion-comparisons) | - | [Suvaditya Mukherjee](https://github.com/suvadityamuk) |
@@ -55,7 +57,7 @@ guided_pipeline = DiffusionPipeline.from_pretrained(
custom_pipeline="clip_guided_stable_diffusion",
clip_model=clip_model,
feature_extractor=feature_extractor,
revision="fp16",
torch_dtype=torch.float16,
)
guided_pipeline.enable_attention_slicing()
@@ -166,7 +168,7 @@ init_image = download_image("https://raw.githubusercontent.com/CompVis/stable-di
prompt = "A fantasy landscape, trending on artstation"
images = pipe.img2img(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
images = pipe.img2img(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
### Inpainting
@@ -176,7 +178,7 @@ init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
prompt = "a cat sitting on a bench"
images = pipe.inpaint(prompt=prompt, init_image=init_image, mask_image=mask_image, strength=0.75).images
images = pipe.inpaint(prompt=prompt, image=init_image, mask_image=mask_image, strength=0.75).images
```
As shown above this one pipeline can run all both "text-to-image", "image-to-image", and "inpainting" in one pipeline.
@@ -206,7 +208,7 @@ import torch
pipe = DiffusionPipeline.from_pretrained(
'hakurei/waifu-diffusion',
custom_pipeline="lpw_stable_diffusion",
revision="fp16",
torch_dtype=torch.float16
)
pipe=pipe.to("cuda")
@@ -273,7 +275,7 @@ diffuser_pipeline = DiffusionPipeline.from_pretrained(
custom_pipeline="speech_to_image_diffusion",
speech_model=model,
speech_processor=processor,
revision="fp16",
torch_dtype=torch.float16,
)
@@ -331,7 +333,7 @@ import torch
pipe = DiffusionPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
custom_pipeline="wildcard_stable_diffusion",
revision="fp16",
torch_dtype=torch.float16,
)
prompt = "__animal__ sitting on a __object__ wearing a __clothing__"
@@ -411,7 +413,7 @@ pipe = DiffusionPipeline.from_pretrained(
custom_pipeline="imagic_stable_diffusion",
scheduler = DDIMScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", clip_sample=False, set_alpha_to_one=False)
).to(device)
generator = th.Generator("cuda").manual_seed(0)
generator = torch.Generator("cuda").manual_seed(0)
seed = 0
prompt = "A photo of Barack Obama smiling with a big grin"
url = 'https://www.dropbox.com/s/6tlwzr73jd1r9yk/obama.png?dl=1'
@@ -420,18 +422,16 @@ init_image = Image.open(BytesIO(response.content)).convert("RGB")
init_image = init_image.resize((512, 512))
res = pipe.train(
prompt,
init_image,
guidance_scale=7.5,
num_inference_steps=50,
image=init_image,
generator=generator)
res = pipe(alpha=1)
res = pipe(alpha=1, guidance_scale=7.5, num_inference_steps=50)
os.makedirs("imagic", exist_ok=True)
image = res.images[0]
image.save('./imagic/imagic_image_alpha_1.png')
res = pipe(alpha=1.5)
res = pipe(alpha=1.5, guidance_scale=7.5, num_inference_steps=50)
image = res.images[0]
image.save('./imagic/imagic_image_alpha_1_5.png')
res = pipe(alpha=2)
res = pipe(alpha=2, guidance_scale=7.5, num_inference_steps=50)
image = res.images[0]
image.save('./imagic/imagic_image_alpha_2.png')
```
@@ -567,7 +567,7 @@ diffuser_pipeline = DiffusionPipeline.from_pretrained(
detection_pipeline=language_detection_pipeline,
translation_model=trans_model,
translation_tokenizer=trans_tokenizer,
revision="fp16",
torch_dtype=torch.float16,
)
@@ -615,7 +615,7 @@ mask_image = PIL.Image.open(mask_path).convert("RGB").resize((512, 512))
pipe = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-inpainting",
custom_pipeline="img2img_inpainting",
revision="fp16",
torch_dtype=torch.float16
)
pipe = pipe.to("cuda")
@@ -687,7 +687,7 @@ pipe = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", custom
pipe = pipe.to("cuda")
prompt = "an astronaut riding a horse on mars"
pipe.set_sampler("sample_heun")
pipe.set_scheduler("sample_heun")
generator = torch.Generator(device="cuda").manual_seed(seed)
image = pipe(prompt, generator=generator, num_inference_steps=20).images[0]
@@ -722,10 +722,97 @@ pipe = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", custom
pipe.scheduler = EulerDiscreteScheduler.from_config(pipe.scheduler.config)
pipe = pipe.to("cuda")
pipe.set_sampler("sample_euler")
pipe.set_scheduler("sample_euler")
generator = torch.Generator(device="cuda").manual_seed(seed)
image = pipe(prompt, generator=generator, num_inference_steps=50).images[0]
```
![diffusers_euler](https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/k_diffusion/astronaut_euler_k_diffusion.png)
### Checkpoint Merger Pipeline
Based on the AUTOMATIC1111/webui for checkpoint merging. This is a custom pipeline that merges upto 3 pretrained model checkpoints as long as they are in the HuggingFace model_index.json format.
The checkpoint merging is currently memory intensive as it modifies the weights of a DiffusionPipeline object in place. Expect atleast 13GB RAM Usage on Kaggle GPU kernels and
on colab you might run out of the 12GB memory even while merging two checkpoints.
Usage:-
```python
from diffusers import DiffusionPipeline
#Return a CheckpointMergerPipeline class that allows you to merge checkpoints.
#The checkpoint passed here is ignored. But still pass one of the checkpoints you plan to
#merge for convenience
pipe = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", custom_pipeline="checkpoint_merger")
#There are multiple possible scenarios:
#The pipeline with the merged checkpoints is returned in all the scenarios
#Compatible checkpoints a.k.a matched model_index.json files. Ignores the meta attributes in model_index.json during comparision.( attrs with _ as prefix )
merged_pipe = pipe.merge(["CompVis/stable-diffusion-v1-4","CompVis/stable-diffusion-v1-2"], interp = "sigmoid", alpha = 0.4)
#Incompatible checkpoints in model_index.json but merge might be possible. Use force = True to ignore model_index.json compatibility
merged_pipe_1 = pipe.merge(["CompVis/stable-diffusion-v1-4","hakurei/waifu-diffusion"], force = True, interp = "sigmoid", alpha = 0.4)
#Three checkpoint merging. Only "add_difference" method actually works on all three checkpoints. Using any other options will ignore the 3rd checkpoint.
merged_pipe_2 = pipe.merge(["CompVis/stable-diffusion-v1-4","hakurei/waifu-diffusion","prompthero/openjourney"], force = True, interp = "add_difference", alpha = 0.4)
prompt = "An astronaut riding a horse on Mars"
image = merged_pipe(prompt).images[0]
```
Some examples along with the merge details:
1. "CompVis/stable-diffusion-v1-4" + "hakurei/waifu-diffusion" ; Sigmoid interpolation; alpha = 0.8
![Stable plus Waifu Sigmoid 0.8](https://huggingface.co/datasets/NagaSaiAbhinay/CheckpointMergerSamples/resolve/main/stability_v1_4_waifu_sig_0.8.png)
2. "hakurei/waifu-diffusion" + "prompthero/openjourney" ; Inverse Sigmoid interpolation; alpha = 0.8
![Stable plus Waifu Sigmoid 0.8](https://huggingface.co/datasets/NagaSaiAbhinay/CheckpointMergerSamples/resolve/main/waifu_openjourney_inv_sig_0.8.png)
3. "CompVis/stable-diffusion-v1-4" + "hakurei/waifu-diffusion" + "prompthero/openjourney"; Add Difference interpolation; alpha = 0.5
![Stable plus Waifu plus openjourney add_diff 0.5](https://huggingface.co/datasets/NagaSaiAbhinay/CheckpointMergerSamples/resolve/main/stable_waifu_openjourney_add_diff_0.5.png)
### Stable Diffusion Comparisons
This Community Pipeline enables the comparison between the 4 checkpoints that exist for Stable Diffusion. They can be found through the following links:
1. [Stable Diffusion v1.1](https://huggingface.co/CompVis/stable-diffusion-v1-1)
2. [Stable Diffusion v1.2](https://huggingface.co/CompVis/stable-diffusion-v1-2)
3. [Stable Diffusion v1.3](https://huggingface.co/CompVis/stable-diffusion-v1-3)
4. [Stable Diffusion v1.4](https://huggingface.co/CompVis/stable-diffusion-v1-4)
```python
from diffusers import DiffusionPipeline
import matplotlib.pyplot as plt
pipe = DiffusionPipeline.from_pretrained('CompVis/stable-diffusion-v1-4', custom_pipeline='suvadityamuk/StableDiffusionComparison')
pipe.enable_attention_slicing()
pipe = pipe.to('cuda')
prompt = "an astronaut riding a horse on mars"
output = pipe(prompt)
plt.subplots(2,2,1)
plt.imshow(output.images[0])
plt.title('Stable Diffusion v1.1')
plt.axis('off')
plt.subplots(2,2,2)
plt.imshow(output.images[1])
plt.title('Stable Diffusion v1.2')
plt.axis('off')
plt.subplots(2,2,3)
plt.imshow(output.images[2])
plt.title('Stable Diffusion v1.3')
plt.axis('off')
plt.subplots(2,2,4)
plt.imshow(output.images[3])
plt.title('Stable Diffusion v1.4')
plt.axis('off')
plt.show()
```python
As a result, you can look at a grid of all 4 generated images being shown together, that captures a difference the advancement of the training between the 4 checkpoints.

View File

@@ -0,0 +1,262 @@
import glob
import os
from typing import Dict, List, Union
import torch
from diffusers import DiffusionPipeline, __version__
from diffusers.pipeline_utils import (
CONFIG_NAME,
DIFFUSERS_CACHE,
ONNX_WEIGHTS_NAME,
SCHEDULER_CONFIG_NAME,
WEIGHTS_NAME,
)
from huggingface_hub import snapshot_download
class CheckpointMergerPipeline(DiffusionPipeline):
"""
A class that that supports merging diffusion models based on the discussion here:
https://github.com/huggingface/diffusers/issues/877
Example usage:-
pipe = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", custom_pipeline="checkpoint_merger.py")
merged_pipe = pipe.merge(["CompVis/stable-diffusion-v1-4","prompthero/openjourney"], interp = 'inv_sigmoid', alpha = 0.8, force = True)
merged_pipe.to('cuda')
prompt = "An astronaut riding a unicycle on Mars"
results = merged_pipe(prompt)
## For more details, see the docstring for the merge method.
"""
def __init__(self):
super().__init__()
def _compare_model_configs(self, dict0, dict1):
if dict0 == dict1:
return True
else:
config0, meta_keys0 = self._remove_meta_keys(dict0)
config1, meta_keys1 = self._remove_meta_keys(dict1)
if config0 == config1:
print(f"Warning !: Mismatch in keys {meta_keys0} and {meta_keys1}.")
return True
return False
def _remove_meta_keys(self, config_dict: Dict):
meta_keys = []
temp_dict = config_dict.copy()
for key in config_dict.keys():
if key.startswith("_"):
temp_dict.pop(key)
meta_keys.append(key)
return (temp_dict, meta_keys)
@torch.no_grad()
def merge(self, pretrained_model_name_or_path_list: List[Union[str, os.PathLike]], **kwargs):
"""
Returns a new pipeline object of the class 'DiffusionPipeline' with the merged checkpoints(weights) of the models passed
in the argument 'pretrained_model_name_or_path_list' as a list.
Parameters:
-----------
pretrained_model_name_or_path_list : A list of valid pretrained model names in the HuggingFace hub or paths to locally stored models in the HuggingFace format.
**kwargs:
Supports all the default DiffusionPipeline.get_config_dict kwargs viz..
cache_dir, resume_download, force_download, proxies, local_files_only, use_auth_token, revision, torch_dtype, device_map.
alpha - The interpolation parameter. Ranges from 0 to 1. It affects the ratio in which the checkpoints are merged. A 0.8 alpha
would mean that the first model checkpoints would affect the final result far less than an alpha of 0.2
interp - The interpolation method to use for the merging. Supports "sigmoid", "inv_sigmoid", "add_difference" and None.
Passing None uses the default interpolation which is weighted sum interpolation. For merging three checkpoints, only "add_difference" is supported.
force - Whether to ignore mismatch in model_config.json for the current models. Defaults to False.
"""
# Default kwargs from DiffusionPipeline
cache_dir = kwargs.pop("cache_dir", DIFFUSERS_CACHE)
resume_download = kwargs.pop("resume_download", False)
force_download = kwargs.pop("force_download", False)
proxies = kwargs.pop("proxies", None)
local_files_only = kwargs.pop("local_files_only", False)
use_auth_token = kwargs.pop("use_auth_token", None)
revision = kwargs.pop("revision", None)
torch_dtype = kwargs.pop("torch_dtype", None)
device_map = kwargs.pop("device_map", None)
alpha = kwargs.pop("alpha", 0.5)
interp = kwargs.pop("interp", None)
print("Recieved list", pretrained_model_name_or_path_list)
checkpoint_count = len(pretrained_model_name_or_path_list)
# Ignore result from model_index_json comparision of the two checkpoints
force = kwargs.pop("force", False)
# If less than 2 checkpoints, nothing to merge. If more than 3, not supported for now.
if checkpoint_count > 3 or checkpoint_count < 2:
raise ValueError(
"Received incorrect number of checkpoints to merge. Ensure that either 2 or 3 checkpoints are being"
" passed."
)
print("Received the right number of checkpoints")
# chkpt0, chkpt1 = pretrained_model_name_or_path_list[0:2]
# chkpt2 = pretrained_model_name_or_path_list[2] if checkpoint_count == 3 else None
# Validate that the checkpoints can be merged
# Step 1: Load the model config and compare the checkpoints. We'll compare the model_index.json first while ignoring the keys starting with '_'
config_dicts = []
for pretrained_model_name_or_path in pretrained_model_name_or_path_list:
if not os.path.isdir(pretrained_model_name_or_path):
config_dict = DiffusionPipeline.get_config_dict(
pretrained_model_name_or_path,
cache_dir=cache_dir,
resume_download=resume_download,
force_download=force_download,
proxies=proxies,
local_files_only=local_files_only,
use_auth_token=use_auth_token,
revision=revision,
)
config_dicts.append(config_dict)
comparison_result = True
for idx in range(1, len(config_dicts)):
comparison_result &= self._compare_model_configs(config_dicts[idx - 1], config_dicts[idx])
if not force and comparison_result is False:
raise ValueError("Incompatible checkpoints. Please check model_index.json for the models.")
print(config_dicts[0], config_dicts[1])
print("Compatible model_index.json files found")
# Step 2: Basic Validation has succeeded. Let's download the models and save them into our local files.
cached_folders = []
for pretrained_model_name_or_path, config_dict in zip(pretrained_model_name_or_path_list, config_dicts):
folder_names = [k for k in config_dict.keys() if not k.startswith("_")]
allow_patterns = [os.path.join(k, "*") for k in folder_names]
allow_patterns += [
WEIGHTS_NAME,
SCHEDULER_CONFIG_NAME,
CONFIG_NAME,
ONNX_WEIGHTS_NAME,
DiffusionPipeline.config_name,
]
requested_pipeline_class = config_dict.get("_class_name")
user_agent = {"diffusers": __version__, "pipeline_class": requested_pipeline_class}
cached_folder = snapshot_download(
pretrained_model_name_or_path,
cache_dir=cache_dir,
resume_download=resume_download,
proxies=proxies,
local_files_only=local_files_only,
use_auth_token=use_auth_token,
revision=revision,
allow_patterns=allow_patterns,
user_agent=user_agent,
)
print("Cached Folder", cached_folder)
cached_folders.append(cached_folder)
# Step 3:-
# Load the first checkpoint as a diffusion pipeline and modify it's module state_dict in place
final_pipe = DiffusionPipeline.from_pretrained(
cached_folders[0], torch_dtype=torch_dtype, device_map=device_map
)
checkpoint_path_2 = None
if len(cached_folders) > 2:
checkpoint_path_2 = os.path.join(cached_folders[2])
if interp == "sigmoid":
theta_func = CheckpointMergerPipeline.sigmoid
elif interp == "inv_sigmoid":
theta_func = CheckpointMergerPipeline.inv_sigmoid
elif interp == "add_diff":
theta_func = CheckpointMergerPipeline.add_difference
else:
theta_func = CheckpointMergerPipeline.weighted_sum
# Find each module's state dict.
for attr in final_pipe.config.keys():
if not attr.startswith("_"):
checkpoint_path_1 = os.path.join(cached_folders[1], attr)
if os.path.exists(checkpoint_path_1):
files = glob.glob(os.path.join(checkpoint_path_1, "*.bin"))
checkpoint_path_1 = files[0] if len(files) > 0 else None
if checkpoint_path_2 is not None and os.path.exists(checkpoint_path_2):
files = glob.glob(os.path.join(checkpoint_path_2, "*.bin"))
checkpoint_path_2 = files[0] if len(files) > 0 else None
# For an attr if both checkpoint_path_1 and 2 are None, ignore.
# If atleast one is present, deal with it according to interp method, of course only if the state_dict keys match.
if checkpoint_path_1 is None and checkpoint_path_2 is None:
print("SKIPPING ATTR ", attr)
continue
try:
module = getattr(final_pipe, attr)
theta_0 = getattr(module, "state_dict")
theta_0 = theta_0()
update_theta_0 = getattr(module, "load_state_dict")
theta_1 = torch.load(checkpoint_path_1)
theta_2 = torch.load(checkpoint_path_2) if checkpoint_path_2 else None
if not theta_0.keys() == theta_1.keys():
print("SKIPPING ATTR ", attr, " DUE TO MISMATCH")
continue
if theta_2 and not theta_1.keys() == theta_2.keys():
print("SKIPPING ATTR ", attr, " DUE TO MISMATCH")
except:
print("SKIPPING ATTR ", attr)
continue
print("Found dicts for")
print(attr)
print(checkpoint_path_1)
print(checkpoint_path_2)
for key in theta_0.keys():
if theta_2:
theta_0[key] = theta_func(theta_0[key], theta_1[key], theta_2[key], alpha)
else:
theta_0[key] = theta_func(theta_0[key], theta_1[key], None, alpha)
del theta_1
del theta_2
update_theta_0(theta_0)
del theta_0
print("Diffusion pipeline successfully updated with merged weights")
return final_pipe
@staticmethod
def weighted_sum(theta0, theta1, theta2, alpha):
return ((1 - alpha) * theta0) + (alpha * theta1)
# Smoothstep (https://en.wikipedia.org/wiki/Smoothstep)
@staticmethod
def sigmoid(theta0, theta1, theta2, alpha):
alpha = alpha * alpha * (3 - (2 * alpha))
return theta0 + ((theta1 - theta0) * alpha)
# Inverse Smoothstep (https://en.wikipedia.org/wiki/Smoothstep)
@staticmethod
def inv_sigmoid(theta0, theta1, theta2, alpha):
import math
alpha = 0.5 - math.sin(math.asin(1.0 - 2.0 * alpha) / 3.0)
return theta0 + ((theta1 - theta0) * alpha)
@staticmethod
def add_difference(theta0, theta1, theta2, alpha):
return theta0 + (theta1 - theta2) * (1.0 - alpha)

View File

@@ -17,7 +17,7 @@ from diffusers.pipeline_utils import DiffusionPipeline
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
from diffusers.pipelines.stable_diffusion.safety_checker import StableDiffusionSafetyChecker
from diffusers.schedulers import DDIMScheduler, LMSDiscreteScheduler, PNDMScheduler
from diffusers.utils import logging
from diffusers.utils import deprecate, logging
# TODO: remove and import from diffusers.utils when the new version of diffusers is released
from packaging import version
@@ -133,7 +133,7 @@ class ImagicStableDiffusionPipeline(DiffusionPipeline):
def train(
self,
prompt: Union[str, List[str]],
init_image: Union[torch.FloatTensor, PIL.Image.Image],
image: Union[torch.FloatTensor, PIL.Image.Image],
height: Optional[int] = 512,
width: Optional[int] = 512,
generator: Optional[torch.Generator] = None,
@@ -184,6 +184,10 @@ class ImagicStableDiffusionPipeline(DiffusionPipeline):
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
(nsfw) content, according to the `safety_checker`.
"""
message = "Please use `image` instead of `init_image`."
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
image = init_image or image
accelerator = Accelerator(
gradient_accumulation_steps=1,
mixed_precision="fp16",
@@ -241,14 +245,14 @@ class ImagicStableDiffusionPipeline(DiffusionPipeline):
lr=embedding_learning_rate,
)
if isinstance(init_image, PIL.Image.Image):
init_image = preprocess(init_image)
if isinstance(image, PIL.Image.Image):
image = preprocess(image)
latents_dtype = text_embeddings.dtype
init_image = init_image.to(device=self.device, dtype=latents_dtype)
init_latent_image_dist = self.vae.encode(init_image).latent_dist
init_image_latents = init_latent_image_dist.sample(generator=generator)
init_image_latents = 0.18215 * init_image_latents
image = image.to(device=self.device, dtype=latents_dtype)
init_latent_image_dist = self.vae.encode(image).latent_dist
image_latents = init_latent_image_dist.sample(generator=generator)
image_latents = 0.18215 * image_latents
progress_bar = tqdm(range(text_embedding_optimization_steps), disable=not accelerator.is_local_main_process)
progress_bar.set_description("Steps")
@@ -259,12 +263,12 @@ class ImagicStableDiffusionPipeline(DiffusionPipeline):
for _ in range(text_embedding_optimization_steps):
with accelerator.accumulate(text_embeddings):
# Sample noise that we'll add to the latents
noise = torch.randn(init_image_latents.shape).to(init_image_latents.device)
timesteps = torch.randint(1000, (1,), device=init_image_latents.device)
noise = torch.randn(image_latents.shape).to(image_latents.device)
timesteps = torch.randint(1000, (1,), device=image_latents.device)
# Add noise to the latents according to the noise magnitude at each timestep
# (this is the forward diffusion process)
noisy_latents = self.scheduler.add_noise(init_image_latents, noise, timesteps)
noisy_latents = self.scheduler.add_noise(image_latents, noise, timesteps)
# Predict the noise residual
noise_pred = self.unet(noisy_latents, timesteps, text_embeddings).sample
@@ -301,12 +305,12 @@ class ImagicStableDiffusionPipeline(DiffusionPipeline):
for _ in range(model_fine_tuning_optimization_steps):
with accelerator.accumulate(self.unet.parameters()):
# Sample noise that we'll add to the latents
noise = torch.randn(init_image_latents.shape).to(init_image_latents.device)
timesteps = torch.randint(1000, (1,), device=init_image_latents.device)
noise = torch.randn(image_latents.shape).to(image_latents.device)
timesteps = torch.randint(1000, (1,), device=image_latents.device)
# Add noise to the latents according to the noise magnitude at each timestep
# (this is the forward diffusion process)
noisy_latents = self.scheduler.add_noise(init_image_latents, noise, timesteps)
noisy_latents = self.scheduler.add_noise(image_latents, noise, timesteps)
# Predict the noise residual
noise_pred = self.unet(noisy_latents, timesteps, text_embeddings).sample

View File

@@ -5,39 +5,37 @@ from typing import Callable, List, Optional, Union
import numpy as np
import torch
import diffusers
import PIL
from diffusers.configuration_utils import FrozenDict
from diffusers import SchedulerMixin, StableDiffusionPipeline
from diffusers.models import AutoencoderKL, UNet2DConditionModel
from diffusers.pipeline_utils import DiffusionPipeline
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
from diffusers.pipelines.stable_diffusion.safety_checker import StableDiffusionSafetyChecker
from diffusers.schedulers import DDIMScheduler, LMSDiscreteScheduler, PNDMScheduler
from diffusers.utils import deprecate, is_accelerate_available, logging
# TODO: remove and import from diffusers.utils when the new version of diffusers is released
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput, StableDiffusionSafetyChecker
from diffusers.utils import deprecate, logging
from packaging import version
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
if version.parse(version.parse(PIL.__version__).base_version) >= version.parse("9.1.0"):
PIL_INTERPOLATION = {
"linear": PIL.Image.Resampling.BILINEAR,
"bilinear": PIL.Image.Resampling.BILINEAR,
"bicubic": PIL.Image.Resampling.BICUBIC,
"lanczos": PIL.Image.Resampling.LANCZOS,
"nearest": PIL.Image.Resampling.NEAREST,
}
else:
PIL_INTERPOLATION = {
"linear": PIL.Image.LINEAR,
"bilinear": PIL.Image.BILINEAR,
"bicubic": PIL.Image.BICUBIC,
"lanczos": PIL.Image.LANCZOS,
"nearest": PIL.Image.NEAREST,
}
try:
from diffusers.utils import PIL_INTERPOLATION
except ImportError:
if version.parse(version.parse(PIL.__version__).base_version) >= version.parse("9.1.0"):
PIL_INTERPOLATION = {
"linear": PIL.Image.Resampling.BILINEAR,
"bilinear": PIL.Image.Resampling.BILINEAR,
"bicubic": PIL.Image.Resampling.BICUBIC,
"lanczos": PIL.Image.Resampling.LANCZOS,
"nearest": PIL.Image.Resampling.NEAREST,
}
else:
PIL_INTERPOLATION = {
"linear": PIL.Image.LINEAR,
"bilinear": PIL.Image.BILINEAR,
"bicubic": PIL.Image.BICUBIC,
"lanczos": PIL.Image.LANCZOS,
"nearest": PIL.Image.NEAREST,
}
# ------------------------------------------------------------------------------
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
re_attention = re.compile(
@@ -146,7 +144,7 @@ def parse_prompt_attention(text):
return res
def get_prompts_with_weights(pipe: DiffusionPipeline, prompt: List[str], max_length: int):
def get_prompts_with_weights(pipe: StableDiffusionPipeline, prompt: List[str], max_length: int):
r"""
Tokenize a list of prompts and return its tokens with weights of each token.
@@ -207,7 +205,7 @@ def pad_tokens_and_weights(tokens, weights, max_length, bos, eos, no_boseos_midd
def get_unweighted_text_embeddings(
pipe: DiffusionPipeline,
pipe: StableDiffusionPipeline,
text_input: torch.Tensor,
chunk_length: int,
no_boseos_middle: Optional[bool] = True,
@@ -247,10 +245,10 @@ def get_unweighted_text_embeddings(
def get_weighted_text_embeddings(
pipe: DiffusionPipeline,
pipe: StableDiffusionPipeline,
prompt: Union[str, List[str]],
uncond_prompt: Optional[Union[str, List[str]]] = None,
max_embeddings_multiples: Optional[int] = 1,
max_embeddings_multiples: Optional[int] = 3,
no_boseos_middle: Optional[bool] = False,
skip_parsing: Optional[bool] = False,
skip_weighting: Optional[bool] = False,
@@ -264,14 +262,14 @@ def get_weighted_text_embeddings(
Also, to regularize of the embedding, the weighted embedding would be scaled to preserve the original mean.
Args:
pipe (`DiffusionPipeline`):
pipe (`StableDiffusionPipeline`):
Pipe to provide access to the tokenizer and the text encoder.
prompt (`str` or `List[str]`):
The prompt or prompts to guide the image generation.
uncond_prompt (`str` or `List[str]`):
The unconditional prompt or prompts for guide the image generation. If unconditional prompt
is provided, the embeddings of prompt and uncond_prompt are concatenated.
max_embeddings_multiples (`int`, *optional*, defaults to `1`):
max_embeddings_multiples (`int`, *optional*, defaults to `3`):
The max multiple length of prompt embeddings compared to the max output length of text encoder.
no_boseos_middle (`bool`, *optional*, defaults to `False`):
If the length of text token is multiples of the capacity of text encoder, whether reserve the starting and
@@ -387,11 +385,11 @@ def preprocess_image(image):
return 2.0 * image - 1.0
def preprocess_mask(mask):
def preprocess_mask(mask, scale_factor=8):
mask = mask.convert("L")
w, h = mask.size
w, h = map(lambda x: x - x % 32, (w, h)) # resize to integer multiple of 32
mask = mask.resize((w // 8, h // 8), resample=PIL_INTERPOLATION["nearest"])
mask = mask.resize((w // scale_factor, h // scale_factor), resample=PIL_INTERPOLATION["nearest"])
mask = np.array(mask).astype(np.float32) / 255.0
mask = np.tile(mask, (4, 1, 1))
mask = mask[None].transpose(0, 1, 2, 3) # what does this step do?
@@ -400,7 +398,7 @@ def preprocess_mask(mask):
return mask
class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
class StableDiffusionLongPromptWeightingPipeline(StableDiffusionPipeline):
r"""
Pipeline for text-to-image generation using Stable Diffusion without tokens length limit, and support parsing
weighting in prompt.
@@ -429,133 +427,245 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
def __init__(
if version.parse(version.parse(diffusers.__version__).base_version) >= version.parse("0.9.0"):
def __init__(
self,
vae: AutoencoderKL,
text_encoder: CLIPTextModel,
tokenizer: CLIPTokenizer,
unet: UNet2DConditionModel,
scheduler: SchedulerMixin,
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPFeatureExtractor,
requires_safety_checker: bool = True,
):
super().__init__(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
requires_safety_checker=requires_safety_checker,
)
self.__init__additional__()
else:
def __init__(
self,
vae: AutoencoderKL,
text_encoder: CLIPTextModel,
tokenizer: CLIPTokenizer,
unet: UNet2DConditionModel,
scheduler: SchedulerMixin,
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPFeatureExtractor,
):
super().__init__(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
)
self.__init__additional__()
def __init__additional__(self):
if not hasattr(self, "vae_scale_factor"):
setattr(self, "vae_scale_factor", 2 ** (len(self.vae.config.block_out_channels) - 1))
@property
def _execution_device(self):
r"""
Returns the device on which the pipeline's models will be executed. After calling
`pipeline.enable_sequential_cpu_offload()` the execution device can only be inferred from Accelerate's module
hooks.
"""
if self.device != torch.device("meta") or not hasattr(self.unet, "_hf_hook"):
return self.device
for module in self.unet.modules():
if (
hasattr(module, "_hf_hook")
and hasattr(module._hf_hook, "execution_device")
and module._hf_hook.execution_device is not None
):
return torch.device(module._hf_hook.execution_device)
return self.device
def _encode_prompt(
self,
vae: AutoencoderKL,
text_encoder: CLIPTextModel,
tokenizer: CLIPTokenizer,
unet: UNet2DConditionModel,
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPFeatureExtractor,
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt,
max_embeddings_multiples,
):
super().__init__()
if hasattr(scheduler.config, "steps_offset") and scheduler.config.steps_offset != 1:
deprecation_message = (
f"The configuration file of this scheduler: {scheduler} is outdated. `steps_offset`"
f" should be set to 1 instead of {scheduler.config.steps_offset}. Please make sure "
"to update the config accordingly as leaving `steps_offset` might led to incorrect results"
" in future versions. If you have downloaded this checkpoint from the Hugging Face Hub,"
" it would be very nice if you could open a Pull request for the `scheduler/scheduler_config.json`"
" file"
)
deprecate("steps_offset!=1", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(scheduler.config)
new_config["steps_offset"] = 1
scheduler._internal_dict = FrozenDict(new_config)
if hasattr(scheduler.config, "clip_sample") and scheduler.config.clip_sample is True:
deprecation_message = (
f"The configuration file of this scheduler: {scheduler} has not set the configuration `clip_sample`."
" `clip_sample` should be set to False in the configuration file. Please make sure to update the"
" config accordingly as not setting `clip_sample` in the config might lead to incorrect results in"
" future versions. If you have downloaded this checkpoint from the Hugging Face Hub, it would be very"
" nice if you could open a Pull request for the `scheduler/scheduler_config.json` file"
)
deprecate("clip_sample not set", "1.0.0", deprecation_message, standard_warn=False)
new_config = dict(scheduler.config)
new_config["clip_sample"] = False
scheduler._internal_dict = FrozenDict(new_config)
if safety_checker is None:
logger.warning(
f"You have disabled the safety checker for {self.__class__} by passing `safety_checker=None`. Ensure"
" that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered"
" results in services or applications open to the public. Both the diffusers team and Hugging Face"
" strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling"
" it only for use-cases that involve analyzing network behavior or auditing its results. For more"
" information, please have a look at https://github.com/huggingface/diffusers/pull/254 ."
)
self.register_modules(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
)
def enable_xformers_memory_efficient_attention(self):
r"""
Enable memory efficient attention as implemented in xformers.
When this option is enabled, you should observe lower GPU memory usage and a potential speed up at inference
time. Speed up at training time is not guaranteed.
Warning: When Memory Efficient Attention and Sliced attention are both enabled, the Memory Efficient Attention
is used.
"""
self.unet.set_use_memory_efficient_attention_xformers(True)
def disable_xformers_memory_efficient_attention(self):
r"""
Disable memory efficient attention as implemented in xformers.
"""
self.unet.set_use_memory_efficient_attention_xformers(False)
def enable_attention_slicing(self, slice_size: Optional[Union[str, int]] = "auto"):
r"""
Enable sliced attention computation.
When this option is enabled, the attention module will split the input tensor in slices, to compute attention
in several steps. This is useful to save some memory in exchange for a small speed decrease.
Encodes the prompt into text encoder hidden states.
Args:
slice_size (`str` or `int`, *optional*, defaults to `"auto"`):
When `"auto"`, halves the input to the attention heads, so attention will be computed in two steps. If
a number is provided, uses as many slices as `attention_head_dim // slice_size`. In this case,
`attention_head_dim` must be a multiple of `slice_size`.
prompt (`str` or `list(int)`):
prompt to be encoded
device: (`torch.device`):
torch device
num_images_per_prompt (`int`):
number of images that should be generated per prompt
do_classifier_free_guidance (`bool`):
whether to use classifier free guidance or not
negative_prompt (`str` or `List[str]`):
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
if `guidance_scale` is less than `1`).
max_embeddings_multiples (`int`, *optional*, defaults to `3`):
The max multiple length of prompt embeddings compared to the max output length of text encoder.
"""
if slice_size == "auto":
# half the attention head size is usually a good trade-off between
# speed and memory
slice_size = self.unet.config.attention_head_dim // 2
self.unet.set_attention_slice(slice_size)
batch_size = len(prompt) if isinstance(prompt, list) else 1
def disable_attention_slicing(self):
r"""
Disable sliced attention computation. If `enable_attention_slicing` was previously invoked, this method will go
back to computing attention in one step.
"""
# set slice_size = `None` to disable `attention slicing`
self.enable_attention_slicing(None)
if negative_prompt is None:
negative_prompt = [""] * batch_size
elif isinstance(negative_prompt, str):
negative_prompt = [negative_prompt] * batch_size
if batch_size != len(negative_prompt):
raise ValueError(
f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:"
f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches"
" the batch size of `prompt`."
)
def enable_sequential_cpu_offload(self):
r"""
Offloads all models to CPU using accelerate, significantly reducing memory usage. When called, unet,
text_encoder, vae and safety checker have their state dicts saved to CPU and then are moved to a
`torch.device('meta') and loaded to GPU only when their specific submodule has its `forward` method called.
"""
if is_accelerate_available():
from accelerate import cpu_offload
text_embeddings, uncond_embeddings = get_weighted_text_embeddings(
pipe=self,
prompt=prompt,
uncond_prompt=negative_prompt if do_classifier_free_guidance else None,
max_embeddings_multiples=max_embeddings_multiples,
)
bs_embed, seq_len, _ = text_embeddings.shape
text_embeddings = text_embeddings.repeat(1, num_images_per_prompt, 1)
text_embeddings = text_embeddings.view(bs_embed * num_images_per_prompt, seq_len, -1)
if do_classifier_free_guidance:
bs_embed, seq_len, _ = uncond_embeddings.shape
uncond_embeddings = uncond_embeddings.repeat(1, num_images_per_prompt, 1)
uncond_embeddings = uncond_embeddings.view(bs_embed * num_images_per_prompt, seq_len, -1)
text_embeddings = torch.cat([uncond_embeddings, text_embeddings])
return text_embeddings
def check_inputs(self, prompt, height, width, strength, callback_steps):
if not isinstance(prompt, str) and not isinstance(prompt, list):
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}")
if strength < 0 or strength > 1:
raise ValueError(f"The value of strength should in [0.0, 1.0] but is {strength}")
if height % 8 != 0 or width % 8 != 0:
raise ValueError(f"`height` and `width` have to be divisible by 8 but are {height} and {width}.")
if (callback_steps is None) or (
callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0)
):
raise ValueError(
f"`callback_steps` has to be a positive integer but is {callback_steps} of type"
f" {type(callback_steps)}."
)
def get_timesteps(self, num_inference_steps, strength, device, is_text2img):
if is_text2img:
return self.scheduler.timesteps.to(device), num_inference_steps
else:
raise ImportError("Please install accelerate via `pip install accelerate`")
# get the original timestep using init_timestep
offset = self.scheduler.config.get("steps_offset", 0)
init_timestep = int(num_inference_steps * strength) + offset
init_timestep = min(init_timestep, num_inference_steps)
device = self.device
t_start = max(num_inference_steps - init_timestep + offset, 0)
timesteps = self.scheduler.timesteps[t_start:].to(device)
return timesteps, num_inference_steps - t_start
for cpu_offloaded_model in [self.unet, self.text_encoder, self.vae, self.safety_checker]:
if cpu_offloaded_model is not None:
cpu_offload(cpu_offloaded_model, device)
def run_safety_checker(self, image, device, dtype):
if self.safety_checker is not None:
safety_checker_input = self.feature_extractor(self.numpy_to_pil(image), return_tensors="pt").to(device)
image, has_nsfw_concept = self.safety_checker(
images=image, clip_input=safety_checker_input.pixel_values.to(dtype)
)
else:
has_nsfw_concept = None
return image, has_nsfw_concept
def decode_latents(self, latents):
latents = 1 / 0.18215 * latents
image = self.vae.decode(latents).sample
image = (image / 2 + 0.5).clamp(0, 1)
# we always cast to float32 as this does not cause significant overhead and is compatible with bfloa16
image = image.cpu().permute(0, 2, 3, 1).float().numpy()
return image
def prepare_extra_step_kwargs(self, generator, eta):
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers.
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502
# and should be between [0, 1]
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys())
extra_step_kwargs = {}
if accepts_eta:
extra_step_kwargs["eta"] = eta
# check if the scheduler accepts generator
accepts_generator = "generator" in set(inspect.signature(self.scheduler.step).parameters.keys())
if accepts_generator:
extra_step_kwargs["generator"] = generator
return extra_step_kwargs
def prepare_latents(self, image, timestep, batch_size, height, width, dtype, device, generator, latents=None):
if image is None:
shape = (
batch_size,
self.unet.in_channels,
height // self.vae_scale_factor,
width // self.vae_scale_factor,
)
if latents is None:
if device.type == "mps":
# randn does not work reproducibly on mps
latents = torch.randn(shape, generator=generator, device="cpu", dtype=dtype).to(device)
else:
latents = torch.randn(shape, generator=generator, device=device, dtype=dtype)
else:
if latents.shape != shape:
raise ValueError(f"Unexpected latents shape, got {latents.shape}, expected {shape}")
latents = latents.to(device)
# scale the initial noise by the standard deviation required by the scheduler
latents = latents * self.scheduler.init_noise_sigma
return latents, None, None
else:
init_latent_dist = self.vae.encode(image).latent_dist
init_latents = init_latent_dist.sample(generator=generator)
init_latents = 0.18215 * init_latents
init_latents = torch.cat([init_latents] * batch_size, dim=0)
init_latents_orig = init_latents
shape = init_latents.shape
# add noise to latents using the timesteps
if device.type == "mps":
noise = torch.randn(shape, generator=generator, device="cpu", dtype=dtype).to(device)
else:
noise = torch.randn(shape, generator=generator, device=device, dtype=dtype)
latents = self.scheduler.add_noise(init_latents, noise, timestep)
return latents, init_latents_orig, noise
@torch.no_grad()
def __call__(
self,
prompt: Union[str, List[str]],
negative_prompt: Optional[Union[str, List[str]]] = None,
init_image: Union[torch.FloatTensor, PIL.Image.Image] = None,
image: Union[torch.FloatTensor, PIL.Image.Image] = None,
mask_image: Union[torch.FloatTensor, PIL.Image.Image] = None,
height: int = 512,
width: int = 512,
@@ -583,11 +693,11 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
negative_prompt (`str` or `List[str]`, *optional*):
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
if `guidance_scale` is less than `1`).
init_image (`torch.FloatTensor` or `PIL.Image.Image`):
image (`torch.FloatTensor` or `PIL.Image.Image`):
`Image`, or tensor representing an image batch, that will be used as the starting point for the
process.
mask_image (`torch.FloatTensor` or `PIL.Image.Image`):
`Image`, or tensor representing an image batch, to mask `init_image`. White pixels in the mask will be
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be
replaced by noise and therefore repainted, while black pixels will be preserved. If `mask_image` is a
PIL image, it will be converted to a single channel (luminance) before use. If it's a tensor, it should
contain one color channel (L) instead of 3, so the expected shape would be `(B, H, W, 1)`.
@@ -605,11 +715,11 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`,
usually at the expense of lower image quality.
strength (`float`, *optional*, defaults to 0.8):
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1.
`image` will be used as a starting point, adding more noise to it the larger the `strength`. The
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
noise will be maximum and the denoising process will run for the full number of iterations specified in
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
num_images_per_prompt (`int`, *optional*, defaults to 1):
The number of images to generate per prompt.
eta (`float`, *optional*, defaults to 0.0):
@@ -648,170 +758,71 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
(nsfw) content, according to the `safety_checker`.
"""
message = "Please use `image` instead of `init_image`."
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
image = init_image or image
if isinstance(prompt, str):
batch_size = 1
prompt = [prompt]
elif isinstance(prompt, list):
batch_size = len(prompt)
else:
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}")
# 0. Default height and width to unet
height = height or self.unet.config.sample_size * self.vae_scale_factor
width = width or self.unet.config.sample_size * self.vae_scale_factor
if strength < 0 or strength > 1:
raise ValueError(f"The value of strength should in [0.0, 1.0] but is {strength}")
if height % 8 != 0 or width % 8 != 0:
raise ValueError(f"`height` and `width` have to be divisible by 8 but are {height} and {width}.")
if (callback_steps is None) or (
callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0)
):
raise ValueError(
f"`callback_steps` has to be a positive integer but is {callback_steps} of type"
f" {type(callback_steps)}."
)
# get prompt text embeddings
# 1. Check inputs. Raise error if not correct
self.check_inputs(prompt, height, width, strength, callback_steps)
# 2. Define call parameters
batch_size = 1 if isinstance(prompt, str) else len(prompt)
device = self._execution_device
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
# corresponds to doing no classifier free guidance.
do_classifier_free_guidance = guidance_scale > 1.0
# get unconditional embeddings for classifier free guidance
if negative_prompt is None:
negative_prompt = [""] * batch_size
elif isinstance(negative_prompt, str):
negative_prompt = [negative_prompt] * batch_size
if batch_size != len(negative_prompt):
raise ValueError(
f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:"
f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches"
" the batch size of `prompt`."
)
text_embeddings, uncond_embeddings = get_weighted_text_embeddings(
pipe=self,
prompt=prompt,
uncond_prompt=negative_prompt if do_classifier_free_guidance else None,
max_embeddings_multiples=max_embeddings_multiples,
**kwargs,
# 3. Encode input prompt
text_embeddings = self._encode_prompt(
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt,
max_embeddings_multiples,
)
bs_embed, seq_len, _ = text_embeddings.shape
text_embeddings = text_embeddings.repeat(1, num_images_per_prompt, 1)
text_embeddings = text_embeddings.view(bs_embed * num_images_per_prompt, seq_len, -1)
dtype = text_embeddings.dtype
if do_classifier_free_guidance:
bs_embed, seq_len, _ = uncond_embeddings.shape
uncond_embeddings = uncond_embeddings.repeat(1, num_images_per_prompt, 1)
uncond_embeddings = uncond_embeddings.view(bs_embed * num_images_per_prompt, seq_len, -1)
text_embeddings = torch.cat([uncond_embeddings, text_embeddings])
# set timesteps
self.scheduler.set_timesteps(num_inference_steps)
latents_dtype = text_embeddings.dtype
init_latents_orig = None
mask = None
noise = None
if init_image is None:
# get the initial random noise unless the user supplied it
# Unlike in other pipelines, latents need to be generated in the target device
# for 1-to-1 results reproducibility with the CompVis implementation.
# However this currently doesn't work in `mps`.
latents_shape = (
batch_size * num_images_per_prompt,
self.unet.in_channels,
height // 8,
width // 8,
)
if latents is None:
if self.device.type == "mps":
# randn does not exist on mps
latents = torch.randn(
latents_shape,
generator=generator,
device="cpu",
dtype=latents_dtype,
).to(self.device)
else:
latents = torch.randn(
latents_shape,
generator=generator,
device=self.device,
dtype=latents_dtype,
)
else:
if latents.shape != latents_shape:
raise ValueError(f"Unexpected latents shape, got {latents.shape}, expected {latents_shape}")
latents = latents.to(self.device)
timesteps = self.scheduler.timesteps.to(self.device)
# scale the initial noise by the standard deviation required by the scheduler
latents = latents * self.scheduler.init_noise_sigma
# 4. Preprocess image and mask
if isinstance(image, PIL.Image.Image):
image = preprocess_image(image)
if image is not None:
image = image.to(device=self.device, dtype=dtype)
if isinstance(mask_image, PIL.Image.Image):
mask_image = preprocess_mask(mask_image, self.vae_scale_factor)
if mask_image is not None:
mask = mask_image.to(device=self.device, dtype=dtype)
mask = torch.cat([mask] * batch_size * num_images_per_prompt)
else:
if isinstance(init_image, PIL.Image.Image):
init_image = preprocess_image(init_image)
# encode the init image into latents and scale the latents
init_image = init_image.to(device=self.device, dtype=latents_dtype)
init_latent_dist = self.vae.encode(init_image).latent_dist
init_latents = init_latent_dist.sample(generator=generator)
init_latents = 0.18215 * init_latents
init_latents = torch.cat([init_latents] * batch_size * num_images_per_prompt, dim=0)
init_latents_orig = init_latents
mask = None
# preprocess mask
if mask_image is not None:
if isinstance(mask_image, PIL.Image.Image):
mask_image = preprocess_mask(mask_image)
mask_image = mask_image.to(device=self.device, dtype=latents_dtype)
mask = torch.cat([mask_image] * batch_size * num_images_per_prompt)
# 5. set timesteps
self.scheduler.set_timesteps(num_inference_steps, device=device)
timesteps, num_inference_steps = self.get_timesteps(num_inference_steps, strength, device, image is None)
latent_timestep = timesteps[:1].repeat(batch_size * num_images_per_prompt)
# check sizes
if not mask.shape == init_latents.shape:
raise ValueError("The mask and init_image should be the same size!")
# 6. Prepare latent variables
latents, init_latents_orig, noise = self.prepare_latents(
image,
latent_timestep,
batch_size * num_images_per_prompt,
height,
width,
dtype,
device,
generator,
latents,
)
# get the original timestep using init_timestep
offset = self.scheduler.config.get("steps_offset", 0)
init_timestep = int(num_inference_steps * strength) + offset
init_timestep = min(init_timestep, num_inference_steps)
timesteps = self.scheduler.timesteps[-init_timestep]
timesteps = torch.tensor([timesteps] * batch_size * num_images_per_prompt, device=self.device)
# add noise to latents using the timesteps
if self.device.type == "mps":
# randn does not exist on mps
noise = torch.randn(
init_latents.shape,
generator=generator,
device="cpu",
dtype=latents_dtype,
).to(self.device)
else:
noise = torch.randn(
init_latents.shape,
generator=generator,
device=self.device,
dtype=latents_dtype,
)
latents = self.scheduler.add_noise(init_latents, noise, timesteps)
t_start = max(num_inference_steps - init_timestep + offset, 0)
timesteps = self.scheduler.timesteps[t_start:].to(self.device)
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers.
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502
# and should be between [0, 1]
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys())
extra_step_kwargs = {}
if accepts_eta:
extra_step_kwargs["eta"] = eta
# 7. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
# 8. Denoising loop
for i, t in enumerate(self.progress_bar(timesteps)):
# expand the latents if we are doing classifier free guidance
latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents
@@ -840,30 +851,18 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
if is_cancelled_callback is not None and is_cancelled_callback():
return None
latents = 1 / 0.18215 * latents
image = self.vae.decode(latents).sample
# 9. Post-processing
image = self.decode_latents(latents)
image = (image / 2 + 0.5).clamp(0, 1)
# we always cast to float32 as this does not cause significant overhead and is compatible with bfloa16
image = image.cpu().permute(0, 2, 3, 1).float().numpy()
if self.safety_checker is not None:
safety_checker_input = self.feature_extractor(self.numpy_to_pil(image), return_tensors="pt").to(
self.device
)
image, has_nsfw_concept = self.safety_checker(
images=image,
clip_input=safety_checker_input.pixel_values.to(text_embeddings.dtype),
)
else:
has_nsfw_concept = None
# 10. Run safety checker
image, has_nsfw_concept = self.run_safety_checker(image, device, text_embeddings.dtype)
# 11. Convert to PIL
if output_type == "pil":
image = self.numpy_to_pil(image)
if not return_dict:
return (image, has_nsfw_concept)
return image, has_nsfw_concept
return StableDiffusionPipelineOutput(images=image, nsfw_content_detected=has_nsfw_concept)
@@ -883,6 +882,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
is_cancelled_callback: Optional[Callable[[], bool]] = None,
callback_steps: Optional[int] = 1,
**kwargs,
):
@@ -930,6 +930,9 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
callback (`Callable`, *optional*):
A function that will be called every `callback_steps` steps during inference. The function will be
called with the following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`.
is_cancelled_callback (`Callable`, *optional*):
A function that will be called every `callback_steps` steps during inference. If the function returns
`True`, the inference will be cancelled.
callback_steps (`int`, *optional*, defaults to 1):
The frequency at which the `callback` function will be called. If not specified, the callback will be
called at every step.
@@ -955,13 +958,14 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
output_type=output_type,
return_dict=return_dict,
callback=callback,
is_cancelled_callback=is_cancelled_callback,
callback_steps=callback_steps,
**kwargs,
)
def img2img(
self,
init_image: Union[torch.FloatTensor, PIL.Image.Image],
image: Union[torch.FloatTensor, PIL.Image.Image],
prompt: Union[str, List[str]],
negative_prompt: Optional[Union[str, List[str]]] = None,
strength: float = 0.8,
@@ -974,13 +978,14 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
is_cancelled_callback: Optional[Callable[[], bool]] = None,
callback_steps: Optional[int] = 1,
**kwargs,
):
r"""
Function for image-to-image generation.
Args:
init_image (`torch.FloatTensor` or `PIL.Image.Image`):
image (`torch.FloatTensor` or `PIL.Image.Image`):
`Image`, or tensor representing an image batch, that will be used as the starting point for the
process.
prompt (`str` or `List[str]`):
@@ -989,11 +994,11 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
if `guidance_scale` is less than `1`).
strength (`float`, *optional*, defaults to 0.8):
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1.
`image` will be used as a starting point, adding more noise to it the larger the `strength`. The
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
noise will be maximum and the denoising process will run for the full number of iterations specified in
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
num_inference_steps (`int`, *optional*, defaults to 50):
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
expense of slower inference. This parameter will be modulated by `strength`.
@@ -1022,6 +1027,9 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
callback (`Callable`, *optional*):
A function that will be called every `callback_steps` steps during inference. The function will be
called with the following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`.
is_cancelled_callback (`Callable`, *optional*):
A function that will be called every `callback_steps` steps during inference. If the function returns
`True`, the inference will be cancelled.
callback_steps (`int`, *optional*, defaults to 1):
The frequency at which the `callback` function will be called. If not specified, the callback will be
called at every step.
@@ -1035,7 +1043,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
return self.__call__(
prompt=prompt,
negative_prompt=negative_prompt,
init_image=init_image,
image=image,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
strength=strength,
@@ -1046,13 +1054,14 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
output_type=output_type,
return_dict=return_dict,
callback=callback,
is_cancelled_callback=is_cancelled_callback,
callback_steps=callback_steps,
**kwargs,
)
def inpaint(
self,
init_image: Union[torch.FloatTensor, PIL.Image.Image],
image: Union[torch.FloatTensor, PIL.Image.Image],
mask_image: Union[torch.FloatTensor, PIL.Image.Image],
prompt: Union[str, List[str]],
negative_prompt: Optional[Union[str, List[str]]] = None,
@@ -1066,17 +1075,18 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
is_cancelled_callback: Optional[Callable[[], bool]] = None,
callback_steps: Optional[int] = 1,
**kwargs,
):
r"""
Function for inpaint.
Args:
init_image (`torch.FloatTensor` or `PIL.Image.Image`):
image (`torch.FloatTensor` or `PIL.Image.Image`):
`Image`, or tensor representing an image batch, that will be used as the starting point for the
process. This is the image whose masked region will be inpainted.
mask_image (`torch.FloatTensor` or `PIL.Image.Image`):
`Image`, or tensor representing an image batch, to mask `init_image`. White pixels in the mask will be
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be
replaced by noise and therefore repainted, while black pixels will be preserved. If `mask_image` is a
PIL image, it will be converted to a single channel (luminance) before use. If it's a tensor, it should
contain one color channel (L) instead of 3, so the expected shape would be `(B, H, W, 1)`.
@@ -1088,7 +1098,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
strength (`float`, *optional*, defaults to 0.8):
Conceptually, indicates how much to inpaint the masked area. Must be between 0 and 1. When `strength`
is 1, the denoising process will be run on the masked area for the full number of iterations specified
in `num_inference_steps`. `init_image` will be used as a reference for the masked area, adding more
in `num_inference_steps`. `image` will be used as a reference for the masked area, adding more
noise to that region the larger the `strength`. If `strength` is 0, no inpainting will occur.
num_inference_steps (`int`, *optional*, defaults to 50):
The reference number of denoising steps. More denoising steps usually lead to a higher quality image at
@@ -1118,6 +1128,9 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
callback (`Callable`, *optional*):
A function that will be called every `callback_steps` steps during inference. The function will be
called with the following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`.
is_cancelled_callback (`Callable`, *optional*):
A function that will be called every `callback_steps` steps during inference. If the function returns
`True`, the inference will be cancelled.
callback_steps (`int`, *optional*, defaults to 1):
The frequency at which the `callback` function will be called. If not specified, the callback will be
called at every step.
@@ -1131,7 +1144,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
return self.__call__(
prompt=prompt,
negative_prompt=negative_prompt,
init_image=init_image,
image=image,
mask_image=mask_image,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
@@ -1143,6 +1156,7 @@ class StableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
output_type=output_type,
return_dict=return_dict,
callback=callback,
is_cancelled_callback=is_cancelled_callback,
callback_steps=callback_steps,
**kwargs,
)

View File

@@ -5,34 +5,53 @@ from typing import Callable, List, Optional, Union
import numpy as np
import torch
import diffusers
import PIL
from diffusers import OnnxStableDiffusionPipeline, SchedulerMixin
from diffusers.onnx_utils import OnnxRuntimeModel
from diffusers.pipeline_utils import DiffusionPipeline
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
from diffusers.schedulers import DDIMScheduler, LMSDiscreteScheduler, PNDMScheduler
from diffusers.utils import logging
# TODO: remove and import from diffusers.utils when the new version of diffusers is released
from diffusers.utils import deprecate, logging
from packaging import version
from transformers import CLIPFeatureExtractor, CLIPTokenizer
if version.parse(version.parse(PIL.__version__).base_version) >= version.parse("9.1.0"):
PIL_INTERPOLATION = {
"linear": PIL.Image.Resampling.BILINEAR,
"bilinear": PIL.Image.Resampling.BILINEAR,
"bicubic": PIL.Image.Resampling.BICUBIC,
"lanczos": PIL.Image.Resampling.LANCZOS,
"nearest": PIL.Image.Resampling.NEAREST,
}
else:
PIL_INTERPOLATION = {
"linear": PIL.Image.LINEAR,
"bilinear": PIL.Image.BILINEAR,
"bicubic": PIL.Image.BICUBIC,
"lanczos": PIL.Image.LANCZOS,
"nearest": PIL.Image.NEAREST,
try:
from diffusers.onnx_utils import ORT_TO_NP_TYPE
except ImportError:
ORT_TO_NP_TYPE = {
"tensor(bool)": np.bool_,
"tensor(int8)": np.int8,
"tensor(uint8)": np.uint8,
"tensor(int16)": np.int16,
"tensor(uint16)": np.uint16,
"tensor(int32)": np.int32,
"tensor(uint32)": np.uint32,
"tensor(int64)": np.int64,
"tensor(uint64)": np.uint64,
"tensor(float16)": np.float16,
"tensor(float)": np.float32,
"tensor(double)": np.float64,
}
try:
from diffusers.utils import PIL_INTERPOLATION
except ImportError:
if version.parse(version.parse(PIL.__version__).base_version) >= version.parse("9.1.0"):
PIL_INTERPOLATION = {
"linear": PIL.Image.Resampling.BILINEAR,
"bilinear": PIL.Image.Resampling.BILINEAR,
"bicubic": PIL.Image.Resampling.BICUBIC,
"lanczos": PIL.Image.Resampling.LANCZOS,
"nearest": PIL.Image.Resampling.NEAREST,
}
else:
PIL_INTERPOLATION = {
"linear": PIL.Image.LINEAR,
"bilinear": PIL.Image.BILINEAR,
"bicubic": PIL.Image.BICUBIC,
"lanczos": PIL.Image.LANCZOS,
"nearest": PIL.Image.NEAREST,
}
# ------------------------------------------------------------------------------
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
@@ -262,7 +281,7 @@ def get_weighted_text_embeddings(
Also, to regularize of the embedding, the weighted embedding would be scaled to preserve the original mean.
Args:
pipe (`DiffusionPipeline`):
pipe (`OnnxStableDiffusionPipeline`):
Pipe to provide access to the tokenizer and the text encoder.
prompt (`str` or `List[str]`):
The prompt or prompts to guide the image generation.
@@ -392,11 +411,11 @@ def preprocess_image(image):
return 2.0 * image - 1.0
def preprocess_mask(mask):
def preprocess_mask(mask, scale_factor=8):
mask = mask.convert("L")
w, h = mask.size
w, h = map(lambda x: x - x % 32, (w, h)) # resize to integer multiple of 32
mask = mask.resize((w // 8, h // 8), resample=PIL_INTERPOLATION["nearest"])
mask = mask.resize((w // scale_factor, h // scale_factor), resample=PIL_INTERPOLATION["nearest"])
mask = np.array(mask).astype(np.float32) / 255.0
mask = np.tile(mask, (4, 1, 1))
mask = mask[None].transpose(0, 1, 2, 3) # what does this step do?
@@ -404,7 +423,7 @@ def preprocess_mask(mask):
return mask
class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
class OnnxStableDiffusionLongPromptWeightingPipeline(OnnxStableDiffusionPipeline):
r"""
Pipeline for text-to-image generation using Stable Diffusion without tokens length limit, and support parsing
weighting in prompt.
@@ -412,36 +431,228 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
"""
if version.parse(version.parse(diffusers.__version__).base_version) >= version.parse("0.9.0"):
def __init__(
def __init__(
self,
vae_encoder: OnnxRuntimeModel,
vae_decoder: OnnxRuntimeModel,
text_encoder: OnnxRuntimeModel,
tokenizer: CLIPTokenizer,
unet: OnnxRuntimeModel,
scheduler: SchedulerMixin,
safety_checker: OnnxRuntimeModel,
feature_extractor: CLIPFeatureExtractor,
requires_safety_checker: bool = True,
):
super().__init__(
vae_encoder=vae_encoder,
vae_decoder=vae_decoder,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
requires_safety_checker=requires_safety_checker,
)
self.__init__additional__()
else:
def __init__(
self,
vae_encoder: OnnxRuntimeModel,
vae_decoder: OnnxRuntimeModel,
text_encoder: OnnxRuntimeModel,
tokenizer: CLIPTokenizer,
unet: OnnxRuntimeModel,
scheduler: SchedulerMixin,
safety_checker: OnnxRuntimeModel,
feature_extractor: CLIPFeatureExtractor,
):
super().__init__(
vae_encoder=vae_encoder,
vae_decoder=vae_decoder,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
)
self.__init__additional__()
def __init__additional__(self):
self.unet_in_channels = 4
self.vae_scale_factor = 8
def _encode_prompt(
self,
vae_encoder: OnnxRuntimeModel,
vae_decoder: OnnxRuntimeModel,
text_encoder: OnnxRuntimeModel,
tokenizer: CLIPTokenizer,
unet: OnnxRuntimeModel,
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
safety_checker: OnnxRuntimeModel,
feature_extractor: CLIPFeatureExtractor,
prompt,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt,
max_embeddings_multiples,
):
super().__init__()
self.register_modules(
vae_encoder=vae_encoder,
vae_decoder=vae_decoder,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
r"""
Encodes the prompt into text encoder hidden states.
Args:
prompt (`str` or `list(int)`):
prompt to be encoded
num_images_per_prompt (`int`):
number of images that should be generated per prompt
do_classifier_free_guidance (`bool`):
whether to use classifier free guidance or not
negative_prompt (`str` or `List[str]`):
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
if `guidance_scale` is less than `1`).
max_embeddings_multiples (`int`, *optional*, defaults to `3`):
The max multiple length of prompt embeddings compared to the max output length of text encoder.
"""
batch_size = len(prompt) if isinstance(prompt, list) else 1
if negative_prompt is None:
negative_prompt = [""] * batch_size
elif isinstance(negative_prompt, str):
negative_prompt = [negative_prompt] * batch_size
if batch_size != len(negative_prompt):
raise ValueError(
f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:"
f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches"
" the batch size of `prompt`."
)
text_embeddings, uncond_embeddings = get_weighted_text_embeddings(
pipe=self,
prompt=prompt,
uncond_prompt=negative_prompt if do_classifier_free_guidance else None,
max_embeddings_multiples=max_embeddings_multiples,
)
text_embeddings = text_embeddings.repeat(num_images_per_prompt, 0)
if do_classifier_free_guidance:
uncond_embeddings = uncond_embeddings.repeat(num_images_per_prompt, 0)
text_embeddings = np.concatenate([uncond_embeddings, text_embeddings])
return text_embeddings
def check_inputs(self, prompt, height, width, strength, callback_steps):
if not isinstance(prompt, str) and not isinstance(prompt, list):
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}")
if strength < 0 or strength > 1:
raise ValueError(f"The value of strength should in [0.0, 1.0] but is {strength}")
if height % 8 != 0 or width % 8 != 0:
raise ValueError(f"`height` and `width` have to be divisible by 8 but are {height} and {width}.")
if (callback_steps is None) or (
callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0)
):
raise ValueError(
f"`callback_steps` has to be a positive integer but is {callback_steps} of type"
f" {type(callback_steps)}."
)
def get_timesteps(self, num_inference_steps, strength, is_text2img):
if is_text2img:
return self.scheduler.timesteps, num_inference_steps
else:
# get the original timestep using init_timestep
offset = self.scheduler.config.get("steps_offset", 0)
init_timestep = int(num_inference_steps * strength) + offset
init_timestep = min(init_timestep, num_inference_steps)
t_start = max(num_inference_steps - init_timestep + offset, 0)
timesteps = self.scheduler.timesteps[t_start:]
return timesteps, num_inference_steps - t_start
def run_safety_checker(self, image):
if self.safety_checker is not None:
safety_checker_input = self.feature_extractor(
self.numpy_to_pil(image), return_tensors="np"
).pixel_values.astype(image.dtype)
# There will throw an error if use safety_checker directly and batchsize>1
images, has_nsfw_concept = [], []
for i in range(image.shape[0]):
image_i, has_nsfw_concept_i = self.safety_checker(
clip_input=safety_checker_input[i : i + 1], images=image[i : i + 1]
)
images.append(image_i)
has_nsfw_concept.append(has_nsfw_concept_i[0])
image = np.concatenate(images)
else:
has_nsfw_concept = None
return image, has_nsfw_concept
def decode_latents(self, latents):
latents = 1 / 0.18215 * latents
# image = self.vae_decoder(latent_sample=latents)[0]
# it seems likes there is a strange result for using half-precision vae decoder if batchsize>1
image = np.concatenate(
[self.vae_decoder(latent_sample=latents[i : i + 1])[0] for i in range(latents.shape[0])]
)
image = np.clip(image / 2 + 0.5, 0, 1)
image = image.transpose((0, 2, 3, 1))
return image
def prepare_extra_step_kwargs(self, generator, eta):
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers.
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502
# and should be between [0, 1]
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys())
extra_step_kwargs = {}
if accepts_eta:
extra_step_kwargs["eta"] = eta
# check if the scheduler accepts generator
accepts_generator = "generator" in set(inspect.signature(self.scheduler.step).parameters.keys())
if accepts_generator:
extra_step_kwargs["generator"] = generator
return extra_step_kwargs
def prepare_latents(self, image, timestep, batch_size, height, width, dtype, generator, latents=None):
if image is None:
shape = (
batch_size,
self.unet_in_channels,
height // self.vae_scale_factor,
width // self.vae_scale_factor,
)
if latents is None:
latents = torch.randn(shape, generator=generator, device="cpu").numpy().astype(dtype)
else:
if latents.shape != shape:
raise ValueError(f"Unexpected latents shape, got {latents.shape}, expected {shape}")
# scale the initial noise by the standard deviation required by the scheduler
latents = (torch.from_numpy(latents) * self.scheduler.init_noise_sigma).numpy()
return latents, None, None
else:
init_latents = self.vae_encoder(sample=image)[0]
init_latents = 0.18215 * init_latents
init_latents = np.concatenate([init_latents] * batch_size, axis=0)
init_latents_orig = init_latents
shape = init_latents.shape
# add noise to latents using the timesteps
noise = torch.randn(shape, generator=generator, device="cpu").numpy().astype(dtype)
latents = self.scheduler.add_noise(
torch.from_numpy(init_latents), torch.from_numpy(noise), timestep
).numpy()
return latents, init_latents_orig, noise
@torch.no_grad()
def __call__(
self,
prompt: Union[str, List[str]],
negative_prompt: Optional[Union[str, List[str]]] = None,
init_image: Union[np.ndarray, PIL.Image.Image] = None,
image: Union[np.ndarray, PIL.Image.Image] = None,
mask_image: Union[np.ndarray, PIL.Image.Image] = None,
height: int = 512,
width: int = 512,
@@ -450,7 +661,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
strength: float = 0.8,
num_images_per_prompt: Optional[int] = 1,
eta: float = 0.0,
generator: Optional[np.random.RandomState] = None,
generator: Optional[torch.Generator] = None,
latents: Optional[np.ndarray] = None,
max_embeddings_multiples: Optional[int] = 3,
output_type: Optional[str] = "pil",
@@ -469,11 +680,11 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
negative_prompt (`str` or `List[str]`, *optional*):
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
if `guidance_scale` is less than `1`).
init_image (`np.ndarray` or `PIL.Image.Image`):
image (`np.ndarray` or `PIL.Image.Image`):
`Image`, or tensor representing an image batch, that will be used as the starting point for the
process.
mask_image (`np.ndarray` or `PIL.Image.Image`):
`Image`, or tensor representing an image batch, to mask `init_image`. White pixels in the mask will be
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be
replaced by noise and therefore repainted, while black pixels will be preserved. If `mask_image` is a
PIL image, it will be converted to a single channel (luminance) before use. If it's a tensor, it should
contain one color channel (L) instead of 3, so the expected shape would be `(B, H, W, 1)`.
@@ -491,18 +702,19 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`,
usually at the expense of lower image quality.
strength (`float`, *optional*, defaults to 0.8):
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1.
`image` will be used as a starting point, adding more noise to it the larger the `strength`. The
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
noise will be maximum and the denoising process will run for the full number of iterations specified in
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
num_images_per_prompt (`int`, *optional*, defaults to 1):
The number of images to generate per prompt.
eta (`float`, *optional*, defaults to 0.0):
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[`schedulers.DDIMScheduler`], will be ignored for others.
generator (`np.random.RandomState`, *optional*):
A np.random.RandomState to make generation deterministic.
generator (`torch.Generator`, *optional*):
A [torch generator](https://pytorch.org/docs/stable/generated/torch.Generator.html) to make generation
deterministic.
latents (`np.ndarray`, *optional*):
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
@@ -533,145 +745,82 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
(nsfw) content, according to the `safety_checker`.
"""
message = "Please use `image` instead of `init_image`."
init_image = deprecate("init_image", "0.12.0", message, take_from=kwargs)
image = init_image or image
if isinstance(prompt, str):
batch_size = 1
prompt = [prompt]
elif isinstance(prompt, list):
batch_size = len(prompt)
else:
raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}")
# 0. Default height and width to unet
height = height or self.unet.config.sample_size * self.vae_scale_factor
width = width or self.unet.config.sample_size * self.vae_scale_factor
if strength < 0 or strength > 1:
raise ValueError(f"The value of strength should in [0.0, 1.0] but is {strength}")
if height % 8 != 0 or width % 8 != 0:
raise ValueError(f"`height` and `width` have to be divisible by 8 but are {height} and {width}.")
if (callback_steps is None) or (
callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0)
):
raise ValueError(
f"`callback_steps` has to be a positive integer but is {callback_steps} of type"
f" {type(callback_steps)}."
)
# get prompt text embeddings
# 1. Check inputs. Raise error if not correct
self.check_inputs(prompt, height, width, strength, callback_steps)
# 2. Define call parameters
batch_size = 1 if isinstance(prompt, str) else len(prompt)
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
# corresponds to doing no classifier free guidance.
do_classifier_free_guidance = guidance_scale > 1.0
# get unconditional embeddings for classifier free guidance
if negative_prompt is None:
negative_prompt = [""] * batch_size
elif isinstance(negative_prompt, str):
negative_prompt = [negative_prompt] * batch_size
if batch_size != len(negative_prompt):
raise ValueError(
f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:"
f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches"
" the batch size of `prompt`."
)
if generator is None:
generator = np.random
# 3. Encode input prompt
text_embeddings = self._encode_prompt(
prompt,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt,
max_embeddings_multiples,
)
dtype = text_embeddings.dtype
text_embeddings, uncond_embeddings = get_weighted_text_embeddings(
pipe=self,
prompt=prompt,
uncond_prompt=negative_prompt if do_classifier_free_guidance else None,
max_embeddings_multiples=max_embeddings_multiples,
**kwargs,
# 4. Preprocess image and mask
if isinstance(image, PIL.Image.Image):
image = preprocess_image(image)
if image is not None:
image = image.astype(dtype)
if isinstance(mask_image, PIL.Image.Image):
mask_image = preprocess_mask(mask_image, self.vae_scale_factor)
if mask_image is not None:
mask = mask_image.astype(dtype)
mask = np.concatenate([mask] * batch_size * num_images_per_prompt)
else:
mask = None
# 5. set timesteps
self.scheduler.set_timesteps(num_inference_steps)
timestep_dtype = next(
(input.type for input in self.unet.model.get_inputs() if input.name == "timestep"), "tensor(float)"
)
timestep_dtype = ORT_TO_NP_TYPE[timestep_dtype]
timesteps, num_inference_steps = self.get_timesteps(num_inference_steps, strength, image is None)
latent_timestep = timesteps[:1].repeat(batch_size * num_images_per_prompt)
# 6. Prepare latent variables
latents, init_latents_orig, noise = self.prepare_latents(
image,
latent_timestep,
batch_size * num_images_per_prompt,
height,
width,
dtype,
generator,
latents,
)
text_embeddings = text_embeddings.repeat(num_images_per_prompt, 0)
if do_classifier_free_guidance:
uncond_embeddings = uncond_embeddings.repeat(num_images_per_prompt, 0)
text_embeddings = np.concatenate([uncond_embeddings, text_embeddings])
# set timesteps
self.scheduler.set_timesteps(num_inference_steps)
latents_dtype = text_embeddings.dtype
init_latents_orig = None
mask = None
noise = None
if init_image is None:
latents_shape = (
batch_size * num_images_per_prompt,
4,
height // 8,
width // 8,
)
if latents is None:
latents = generator.randn(*latents_shape).astype(latents_dtype)
elif latents.shape != latents_shape:
raise ValueError(f"Unexpected latents shape, got {latents.shape}, expected {latents_shape}")
timesteps = self.scheduler.timesteps.to(self.device)
# scale the initial noise by the standard deviation required by the scheduler
latents = latents * self.scheduler.init_noise_sigma
else:
if isinstance(init_image, PIL.Image.Image):
init_image = preprocess_image(init_image)
# encode the init image into latents and scale the latents
init_image = init_image.astype(latents_dtype)
init_latents = self.vae_encoder(sample=init_image)[0]
init_latents = 0.18215 * init_latents
init_latents = np.concatenate([init_latents] * batch_size * num_images_per_prompt)
init_latents_orig = init_latents
# preprocess mask
if mask_image is not None:
if isinstance(mask_image, PIL.Image.Image):
mask_image = preprocess_mask(mask_image)
mask_image = mask_image.astype(latents_dtype)
mask = np.concatenate([mask_image] * batch_size * num_images_per_prompt)
# check sizes
if not mask.shape == init_latents.shape:
print(mask.shape, init_latents.shape)
raise ValueError("The mask and init_image should be the same size!")
# get the original timestep using init_timestep
offset = self.scheduler.config.get("steps_offset", 0)
init_timestep = int(num_inference_steps * strength) + offset
init_timestep = min(init_timestep, num_inference_steps)
timesteps = self.scheduler.timesteps[-init_timestep]
timesteps = torch.tensor([timesteps] * batch_size * num_images_per_prompt)
# add noise to latents using the timesteps
noise = generator.randn(*init_latents.shape).astype(latents_dtype)
latents = self.scheduler.add_noise(
torch.from_numpy(init_latents), torch.from_numpy(noise), timesteps
).numpy()
t_start = max(num_inference_steps - init_timestep + offset, 0)
timesteps = self.scheduler.timesteps[t_start:]
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers.
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502
# and should be between [0, 1]
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys())
extra_step_kwargs = {}
if accepts_eta:
extra_step_kwargs["eta"] = eta
# 7. Prepare extra step kwargs. TODO: Logic should ideally just be moved out of the pipeline
extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta)
# 8. Denoising loop
for i, t in enumerate(self.progress_bar(timesteps)):
# expand the latents if we are doing classifier free guidance
latent_model_input = np.concatenate([latents] * 2) if do_classifier_free_guidance else latents
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
latent_model_input = self.scheduler.scale_model_input(torch.from_numpy(latent_model_input), t)
latent_model_input = latent_model_input.numpy()
# predict the noise residual
noise_pred = self.unet(
sample=latent_model_input,
timestep=np.array([t]),
timestep=np.array([t], dtype=timestep_dtype),
encoder_hidden_states=text_embeddings,
)
noise_pred = noise_pred[0]
@@ -682,14 +831,17 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
# compute the previous noisy sample x_t -> x_t-1
latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs).prev_sample.numpy()
scheduler_output = self.scheduler.step(
torch.from_numpy(noise_pred), t, torch.from_numpy(latents), **extra_step_kwargs
)
latents = scheduler_output.prev_sample.numpy()
if mask is not None:
# masking
init_latents_proper = self.scheduler.add_noise(
torch.from_numpy(init_latents_orig),
torch.from_numpy(noise),
torch.tensor([t]),
t,
).numpy()
latents = (init_latents_proper * mask) + (latents * (1 - mask))
@@ -700,38 +852,18 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
if is_cancelled_callback is not None and is_cancelled_callback():
return None
latents = 1 / 0.18215 * latents
# image = self.vae_decoder(latent_sample=latents)[0]
# it seems likes there is a problem for using half-precision vae decoder if batchsize>1
image = []
for i in range(latents.shape[0]):
image.append(self.vae_decoder(latent_sample=latents[i : i + 1])[0])
image = np.concatenate(image)
# 9. Post-processing
image = self.decode_latents(latents)
image = np.clip(image / 2 + 0.5, 0, 1)
image = image.transpose((0, 2, 3, 1))
if self.safety_checker is not None:
safety_checker_input = self.feature_extractor(
self.numpy_to_pil(image), return_tensors="np"
).pixel_values.astype(image.dtype)
# There will throw an error if use safety_checker directly and batchsize>1
images, has_nsfw_concept = [], []
for i in range(image.shape[0]):
image_i, has_nsfw_concept_i = self.safety_checker(
clip_input=safety_checker_input[i : i + 1], images=image[i : i + 1]
)
images.append(image_i)
has_nsfw_concept.append(has_nsfw_concept_i[0])
image = np.concatenate(images)
else:
has_nsfw_concept = None
# 10. Run safety checker
image, has_nsfw_concept = self.run_safety_checker(image)
# 11. Convert to PIL
if output_type == "pil":
image = self.numpy_to_pil(image)
if not return_dict:
return (image, has_nsfw_concept)
return image, has_nsfw_concept
return StableDiffusionPipelineOutput(images=image, nsfw_content_detected=has_nsfw_concept)
@@ -745,7 +877,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
guidance_scale: float = 7.5,
num_images_per_prompt: Optional[int] = 1,
eta: float = 0.0,
generator: Optional[np.random.RandomState] = None,
generator: Optional[torch.Generator] = None,
latents: Optional[np.ndarray] = None,
max_embeddings_multiples: Optional[int] = 3,
output_type: Optional[str] = "pil",
@@ -780,8 +912,9 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
eta (`float`, *optional*, defaults to 0.0):
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[`schedulers.DDIMScheduler`], will be ignored for others.
generator (`np.random.RandomState`, *optional*):
A np.random.RandomState to make generation deterministic.
generator (`torch.Generator`, *optional*):
A [torch generator](https://pytorch.org/docs/stable/generated/torch.Generator.html) to make generation
deterministic.
latents (`np.ndarray`, *optional*):
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
@@ -828,7 +961,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
def img2img(
self,
init_image: Union[np.ndarray, PIL.Image.Image],
image: Union[np.ndarray, PIL.Image.Image],
prompt: Union[str, List[str]],
negative_prompt: Optional[Union[str, List[str]]] = None,
strength: float = 0.8,
@@ -836,7 +969,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
guidance_scale: Optional[float] = 7.5,
num_images_per_prompt: Optional[int] = 1,
eta: Optional[float] = 0.0,
generator: Optional[np.random.RandomState] = None,
generator: Optional[torch.Generator] = None,
max_embeddings_multiples: Optional[int] = 3,
output_type: Optional[str] = "pil",
return_dict: bool = True,
@@ -847,7 +980,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
r"""
Function for image-to-image generation.
Args:
init_image (`np.ndarray` or `PIL.Image.Image`):
image (`np.ndarray` or `PIL.Image.Image`):
`Image`, or ndarray representing an image batch, that will be used as the starting point for the
process.
prompt (`str` or `List[str]`):
@@ -856,11 +989,11 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
The prompt or prompts not to guide the image generation. Ignored when not using guidance (i.e., ignored
if `guidance_scale` is less than `1`).
strength (`float`, *optional*, defaults to 0.8):
Conceptually, indicates how much to transform the reference `init_image`. Must be between 0 and 1.
`init_image` will be used as a starting point, adding more noise to it the larger the `strength`. The
Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1.
`image` will be used as a starting point, adding more noise to it the larger the `strength`. The
number of denoising steps depends on the amount of noise initially added. When `strength` is 1, added
noise will be maximum and the denoising process will run for the full number of iterations specified in
`num_inference_steps`. A value of 1, therefore, essentially ignores `init_image`.
`num_inference_steps`. A value of 1, therefore, essentially ignores `image`.
num_inference_steps (`int`, *optional*, defaults to 50):
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
expense of slower inference. This parameter will be modulated by `strength`.
@@ -875,8 +1008,9 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
eta (`float`, *optional*, defaults to 0.0):
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[`schedulers.DDIMScheduler`], will be ignored for others.
generator (`np.random.RandomState`, *optional*):
A np.random.RandomState to make generation deterministic.
generator (`torch.Generator`, *optional*):
A [torch generator](https://pytorch.org/docs/stable/generated/torch.Generator.html) to make generation
deterministic.
max_embeddings_multiples (`int`, *optional*, defaults to `3`):
The max multiple length of prompt embeddings compared to the max output length of text encoder.
output_type (`str`, *optional*, defaults to `"pil"`):
@@ -901,7 +1035,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
return self.__call__(
prompt=prompt,
negative_prompt=negative_prompt,
init_image=init_image,
image=image,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
strength=strength,
@@ -918,7 +1052,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
def inpaint(
self,
init_image: Union[np.ndarray, PIL.Image.Image],
image: Union[np.ndarray, PIL.Image.Image],
mask_image: Union[np.ndarray, PIL.Image.Image],
prompt: Union[str, List[str]],
negative_prompt: Optional[Union[str, List[str]]] = None,
@@ -927,7 +1061,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
guidance_scale: Optional[float] = 7.5,
num_images_per_prompt: Optional[int] = 1,
eta: Optional[float] = 0.0,
generator: Optional[np.random.RandomState] = None,
generator: Optional[torch.Generator] = None,
max_embeddings_multiples: Optional[int] = 3,
output_type: Optional[str] = "pil",
return_dict: bool = True,
@@ -938,11 +1072,11 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
r"""
Function for inpaint.
Args:
init_image (`np.ndarray` or `PIL.Image.Image`):
image (`np.ndarray` or `PIL.Image.Image`):
`Image`, or tensor representing an image batch, that will be used as the starting point for the
process. This is the image whose masked region will be inpainted.
mask_image (`np.ndarray` or `PIL.Image.Image`):
`Image`, or tensor representing an image batch, to mask `init_image`. White pixels in the mask will be
`Image`, or tensor representing an image batch, to mask `image`. White pixels in the mask will be
replaced by noise and therefore repainted, while black pixels will be preserved. If `mask_image` is a
PIL image, it will be converted to a single channel (luminance) before use. If it's a tensor, it should
contain one color channel (L) instead of 3, so the expected shape would be `(B, H, W, 1)`.
@@ -954,7 +1088,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
strength (`float`, *optional*, defaults to 0.8):
Conceptually, indicates how much to inpaint the masked area. Must be between 0 and 1. When `strength`
is 1, the denoising process will be run on the masked area for the full number of iterations specified
in `num_inference_steps`. `init_image` will be used as a reference for the masked area, adding more
in `num_inference_steps`. `image` will be used as a reference for the masked area, adding more
noise to that region the larger the `strength`. If `strength` is 0, no inpainting will occur.
num_inference_steps (`int`, *optional*, defaults to 50):
The reference number of denoising steps. More denoising steps usually lead to a higher quality image at
@@ -970,8 +1104,9 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
eta (`float`, *optional*, defaults to 0.0):
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[`schedulers.DDIMScheduler`], will be ignored for others.
generator (`np.random.RandomState`, *optional*):
A np.random.RandomState to make generation deterministic.
generator (`torch.Generator`, *optional*):
A [torch generator](https://pytorch.org/docs/stable/generated/torch.Generator.html) to make generation
deterministic.
max_embeddings_multiples (`int`, *optional*, defaults to `3`):
The max multiple length of prompt embeddings compared to the max output length of text encoder.
output_type (`str`, *optional*, defaults to `"pil"`):
@@ -996,7 +1131,7 @@ class OnnxStableDiffusionLongPromptWeightingPipeline(DiffusionPipeline):
return self.__call__(
prompt=prompt,
negative_prompt=negative_prompt,
init_image=init_image,
image=image,
mask_image=mask_image,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,

View File

@@ -19,4 +19,6 @@ class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
model_output = self.unet(image, timestep).sample
scheduler_output = self.scheduler.step(model_output, timestep, image).prev_sample
return scheduler_output
result = scheduler_output - scheduler_output + torch.ones_like(scheduler_output)
return result

View File

@@ -13,6 +13,7 @@
# limitations under the License.
import importlib
import warnings
from typing import Callable, List, Optional, Union
import torch
@@ -21,7 +22,7 @@ from diffusers import LMSDiscreteScheduler
from diffusers.pipeline_utils import DiffusionPipeline
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
from diffusers.utils import is_accelerate_available, logging
from k_diffusion.external import CompVisDenoiser
from k_diffusion.external import CompVisDenoiser, CompVisVDenoiser
logger = logging.get_logger(__name__) # pylint: disable=invalid-name
@@ -33,7 +34,12 @@ class ModelWrapper:
self.alphas_cumprod = alphas_cumprod
def apply_model(self, *args, **kwargs):
return self.model(*args, **kwargs).sample
if len(args) == 3:
encoder_hidden_states = args[-1]
args = args[:2]
if kwargs.get("cond", None) is not None:
encoder_hidden_states = kwargs.pop("cond")
return self.model(*args, encoder_hidden_states=encoder_hidden_states, **kwargs).sample
class StableDiffusionPipeline(DiffusionPipeline):
@@ -63,6 +69,7 @@ class StableDiffusionPipeline(DiffusionPipeline):
feature_extractor ([`CLIPFeatureExtractor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
_optional_components = ["safety_checker", "feature_extractor"]
def __init__(
self,
@@ -99,31 +106,20 @@ class StableDiffusionPipeline(DiffusionPipeline):
)
model = ModelWrapper(unet, scheduler.alphas_cumprod)
self.k_diffusion_model = CompVisDenoiser(model)
if scheduler.prediction_type == "v_prediction":
self.k_diffusion_model = CompVisVDenoiser(model)
else:
self.k_diffusion_model = CompVisDenoiser(model)
def set_sampler(self, scheduler_type: str):
warnings.warn("The `set_sampler` method is deprecated, please use `set_scheduler` instead.")
return self.set_scheduler(scheduler_type)
def set_scheduler(self, scheduler_type: str):
library = importlib.import_module("k_diffusion")
sampling = getattr(library, "sampling")
self.sampler = getattr(sampling, scheduler_type)
def enable_xformers_memory_efficient_attention(self):
r"""
Enable memory efficient attention as implemented in xformers.
When this option is enabled, you should observe lower GPU memory usage and a potential speed up at inference
time. Speed up at training time is not guaranteed.
Warning: When Memory Efficient Attention and Sliced attention are both enabled, the Memory Efficient Attention
is used.
"""
self.unet.set_use_memory_efficient_attention_xformers(True)
def disable_xformers_memory_efficient_attention(self):
r"""
Disable memory efficient attention as implemented in xformers.
"""
self.unet.set_use_memory_efficient_attention_xformers(False)
def enable_attention_slicing(self, slice_size: Optional[Union[str, int]] = "auto"):
r"""
Enable sliced attention computation.
@@ -435,6 +431,7 @@ class StableDiffusionPipeline(DiffusionPipeline):
# 4. Prepare timesteps
self.scheduler.set_timesteps(num_inference_steps, device=text_embeddings.device)
sigmas = self.scheduler.sigmas
sigmas = sigmas.to(text_embeddings.dtype)
# 5. Prepare latent variables
num_channels_latents = self.unet.in_channels
@@ -455,7 +452,7 @@ class StableDiffusionPipeline(DiffusionPipeline):
def model_fn(x, t):
latent_model_input = torch.cat([x] * 2)
noise_pred = self.k_diffusion_model(latent_model_input, t, encoder_hidden_states=text_embeddings)
noise_pred = self.k_diffusion_model(latent_model_input, t, cond=text_embeddings)
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)

View File

@@ -0,0 +1,405 @@
from typing import Any, Callable, Dict, List, Optional, Union
import torch
from diffusers import (
AutoencoderKL,
DDIMScheduler,
DiffusionPipeline,
LMSDiscreteScheduler,
PNDMScheduler,
StableDiffusionPipeline,
UNet2DConditionModel,
)
from diffusers.pipelines.stable_diffusion import StableDiffusionPipelineOutput
from diffusers.pipelines.stable_diffusion.safety_checker import StableDiffusionSafetyChecker
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
pipe1_model_id = "CompVis/stable-diffusion-v1-1"
pipe2_model_id = "CompVis/stable-diffusion-v1-2"
pipe3_model_id = "CompVis/stable-diffusion-v1-3"
pipe4_model_id = "CompVis/stable-diffusion-v1-4"
class StableDiffusionComparisonPipeline(DiffusionPipeline):
r"""
Pipeline for parallel comparison of Stable Diffusion v1-v4
This pipeline inherits from DiffusionPipeline and depends on the use of an Auth Token for
downloading pre-trained checkpoints from Hugging Face Hub.
If using Hugging Face Hub, pass the Model ID for Stable Diffusion v1.4 as the previous 3 checkpoints will be loaded
automatically.
Args:
vae ([`AutoencoderKL`]):
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
text_encoder ([`CLIPTextModel`]):
Frozen text-encoder. Stable Diffusion uses the text portion of
[CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModel), specifically
the [clip-vit-large-patch14](https://huggingface.co/openai/clip-vit-large-patch14) variant.
tokenizer (`CLIPTokenizer`):
Tokenizer of class
[CLIPTokenizer](https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer).
unet ([`UNet2DConditionModel`]): Conditional U-Net architecture to denoise the encoded image latents.
scheduler ([`SchedulerMixin`]):
A scheduler to be used in combination with `unet` to denoise the encoded image latents. Can be one of
[`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`].
safety_checker ([`StableDiffusionMegaSafetyChecker`]):
Classification module that estimates whether generated images could be considered offensive or harmful.
Please, refer to the [model card](https://huggingface.co/runwayml/stable-diffusion-v1-5) for details.
feature_extractor ([`CLIPFeatureExtractor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
def __init__(
self,
vae: AutoencoderKL,
text_encoder: CLIPTextModel,
tokenizer: CLIPTokenizer,
unet: UNet2DConditionModel,
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPFeatureExtractor,
requires_safety_checker: bool = True,
):
super()._init_()
self.pipe1 = StableDiffusionPipeline.from_pretrained(pipe1_model_id)
self.pipe2 = StableDiffusionPipeline.from_pretrained(pipe2_model_id)
self.pipe3 = StableDiffusionPipeline.from_pretrained(pipe3_model_id)
self.pipe4 = StableDiffusionPipeline(
vae=vae,
text_encoder=text_encoder,
tokenizer=tokenizer,
unet=unet,
scheduler=scheduler,
safety_checker=safety_checker,
feature_extractor=feature_extractor,
requires_safety_checker=requires_safety_checker,
)
self.register_modules(pipeline1=self.pipe1, pipeline2=self.pipe2, pipeline3=self.pipe3, pipeline4=self.pipe4)
@property
def layers(self) -> Dict[str, Any]:
return {k: getattr(self, k) for k in self.config.keys() if not k.startswith("_")}
def enable_attention_slicing(self, slice_size: Optional[Union[str, int]] = "auto"):
r"""
Enable sliced attention computation.
When this option is enabled, the attention module will split the input tensor in slices, to compute attention
in several steps. This is useful to save some memory in exchange for a small speed decrease.
Args:
slice_size (`str` or `int`, *optional*, defaults to `"auto"`):
When `"auto"`, halves the input to the attention heads, so attention will be computed in two steps. If
a number is provided, uses as many slices as `attention_head_dim // slice_size`. In this case,
`attention_head_dim` must be a multiple of `slice_size`.
"""
if slice_size == "auto":
# half the attention head size is usually a good trade-off between
# speed and memory
slice_size = self.unet.config.attention_head_dim // 2
self.unet.set_attention_slice(slice_size)
def disable_attention_slicing(self):
r"""
Disable sliced attention computation. If `enable_attention_slicing` was previously invoked, this method will go
back to computing attention in one step.
"""
# set slice_size = `None` to disable `attention slicing`
self.enable_attention_slicing(None)
@torch.no_grad()
def text2img_sd1_1(
self,
prompt: Union[str, List[str]],
height: int = 512,
width: int = 512,
num_inference_steps: int = 50,
guidance_scale: float = 7.5,
negative_prompt: Optional[Union[str, List[str]]] = None,
num_images_per_prompt: Optional[int] = 1,
eta: float = 0.0,
generator: Optional[torch.Generator] = None,
latents: Optional[torch.FloatTensor] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
callback_steps: Optional[int] = 1,
**kwargs,
):
return self.pipe1(
prompt=prompt,
height=height,
width=width,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
negative_prompt=negative_prompt,
num_images_per_prompt=num_images_per_prompt,
eta=eta,
generator=generator,
latents=latents,
output_type=output_type,
return_dict=return_dict,
callback=callback,
callback_steps=callback_steps,
**kwargs,
)
@torch.no_grad()
def text2img_sd1_2(
self,
prompt: Union[str, List[str]],
height: int = 512,
width: int = 512,
num_inference_steps: int = 50,
guidance_scale: float = 7.5,
negative_prompt: Optional[Union[str, List[str]]] = None,
num_images_per_prompt: Optional[int] = 1,
eta: float = 0.0,
generator: Optional[torch.Generator] = None,
latents: Optional[torch.FloatTensor] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
callback_steps: Optional[int] = 1,
**kwargs,
):
return self.pipe2(
prompt=prompt,
height=height,
width=width,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
negative_prompt=negative_prompt,
num_images_per_prompt=num_images_per_prompt,
eta=eta,
generator=generator,
latents=latents,
output_type=output_type,
return_dict=return_dict,
callback=callback,
callback_steps=callback_steps,
**kwargs,
)
@torch.no_grad()
def text2img_sd1_3(
self,
prompt: Union[str, List[str]],
height: int = 512,
width: int = 512,
num_inference_steps: int = 50,
guidance_scale: float = 7.5,
negative_prompt: Optional[Union[str, List[str]]] = None,
num_images_per_prompt: Optional[int] = 1,
eta: float = 0.0,
generator: Optional[torch.Generator] = None,
latents: Optional[torch.FloatTensor] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
callback_steps: Optional[int] = 1,
**kwargs,
):
return self.pipe3(
prompt=prompt,
height=height,
width=width,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
negative_prompt=negative_prompt,
num_images_per_prompt=num_images_per_prompt,
eta=eta,
generator=generator,
latents=latents,
output_type=output_type,
return_dict=return_dict,
callback=callback,
callback_steps=callback_steps,
**kwargs,
)
@torch.no_grad()
def text2img_sd1_4(
self,
prompt: Union[str, List[str]],
height: int = 512,
width: int = 512,
num_inference_steps: int = 50,
guidance_scale: float = 7.5,
negative_prompt: Optional[Union[str, List[str]]] = None,
num_images_per_prompt: Optional[int] = 1,
eta: float = 0.0,
generator: Optional[torch.Generator] = None,
latents: Optional[torch.FloatTensor] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
callback_steps: Optional[int] = 1,
**kwargs,
):
return self.pipe4(
prompt=prompt,
height=height,
width=width,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
negative_prompt=negative_prompt,
num_images_per_prompt=num_images_per_prompt,
eta=eta,
generator=generator,
latents=latents,
output_type=output_type,
return_dict=return_dict,
callback=callback,
callback_steps=callback_steps,
**kwargs,
)
@torch.no_grad()
def _call_(
self,
prompt: Union[str, List[str]],
height: int = 512,
width: int = 512,
num_inference_steps: int = 50,
guidance_scale: float = 7.5,
negative_prompt: Optional[Union[str, List[str]]] = None,
num_images_per_prompt: Optional[int] = 1,
eta: float = 0.0,
generator: Optional[torch.Generator] = None,
latents: Optional[torch.FloatTensor] = None,
output_type: Optional[str] = "pil",
return_dict: bool = True,
callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None,
callback_steps: Optional[int] = 1,
**kwargs,
):
r"""
Function invoked when calling the pipeline for generation. This function will generate 4 results as part
of running all the 4 pipelines for SD1.1-1.4 together in a serial-processing, parallel-invocation fashion.
Args:
prompt (`str` or `List[str]`):
The prompt or prompts to guide the image generation.
height (`int`, optional, defaults to 512):
The height in pixels of the generated image.
width (`int`, optional, defaults to 512):
The width in pixels of the generated image.
num_inference_steps (`int`, optional, defaults to 50):
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
expense of slower inference.
guidance_scale (`float`, optional, defaults to 7.5):
Guidance scale as defined in [Classifier-Free Diffusion Guidance](https://arxiv.org/abs/2207.12598).
`guidance_scale` is defined as `w` of equation 2. of [Imagen
Paper](https://arxiv.org/pdf/2205.11487.pdf). Guidance scale is enabled by setting `guidance_scale >
1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`,
usually at the expense of lower image quality.
eta (`float`, optional, defaults to 0.0):
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[`schedulers.DDIMScheduler`], will be ignored for others.
generator (`torch.Generator`, optional):
A [torch generator](https://pytorch.org/docs/stable/generated/torch.Generator.html) to make generation
deterministic.
latents (`torch.FloatTensor`, optional):
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor will ge generated by sampling using the supplied random `generator`.
output_type (`str`, optional, defaults to `"pil"`):
The output format of the generate image. Choose between
[PIL](https://pillow.readthedocs.io/en/stable/): `PIL.Image.Image` or `np.array`.
return_dict (`bool`, optional, defaults to `True`):
Whether or not to return a [`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] instead of a
plain tuple.
Returns:
[`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] or `tuple`:
[`~pipelines.stable_diffusion.StableDiffusionPipelineOutput`] if `return_dict` is True, otherwise a `tuple.
When returning a tuple, the first element is a list with the generated images, and the second element is a
list of `bool`s denoting whether the corresponding generated image likely represents "not-safe-for-work"
(nsfw) content, according to the `safety_checker`.
"""
device = "cuda" if torch.cuda.is_available() else "cpu"
self.to(device)
# Checks if the height and width are divisible by 8 or not
if height % 8 != 0 or width % 8 != 0:
raise ValueError(f"`height` and `width` must be divisible by 8 but are {height} and {width}.")
# Get first result from Stable Diffusion Checkpoint v1.1
res1 = self.text2img_sd1_1(
prompt=prompt,
height=height,
width=width,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
negative_prompt=negative_prompt,
num_images_per_prompt=num_images_per_prompt,
eta=eta,
generator=generator,
latents=latents,
output_type=output_type,
return_dict=return_dict,
callback=callback,
callback_steps=callback_steps,
**kwargs,
)
# Get first result from Stable Diffusion Checkpoint v1.2
res2 = self.text2img_sd1_2(
prompt=prompt,
height=height,
width=width,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
negative_prompt=negative_prompt,
num_images_per_prompt=num_images_per_prompt,
eta=eta,
generator=generator,
latents=latents,
output_type=output_type,
return_dict=return_dict,
callback=callback,
callback_steps=callback_steps,
**kwargs,
)
# Get first result from Stable Diffusion Checkpoint v1.3
res3 = self.text2img_sd1_3(
prompt=prompt,
height=height,
width=width,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
negative_prompt=negative_prompt,
num_images_per_prompt=num_images_per_prompt,
eta=eta,
generator=generator,
latents=latents,
output_type=output_type,
return_dict=return_dict,
callback=callback,
callback_steps=callback_steps,
**kwargs,
)
# Get first result from Stable Diffusion Checkpoint v1.4
res4 = self.text2img_sd1_4(
prompt=prompt,
height=height,
width=width,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,
negative_prompt=negative_prompt,
num_images_per_prompt=num_images_per_prompt,
eta=eta,
generator=generator,
latents=latents,
output_type=output_type,
return_dict=return_dict,
callback=callback,
callback_steps=callback_steps,
**kwargs,
)
# Get all result images into a single list and pass it via StableDiffusionPipelineOutput for final result
return StableDiffusionPipelineOutput([res1[0], res2[0], res3[0], res4[0]])

View File

@@ -50,6 +50,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
feature_extractor ([`CLIPFeatureExtractor`]):
Model that extracts features from generated images to be used as inputs for the `safety_checker`.
"""
_optional_components = ["safety_checker", "feature_extractor"]
def __init__(
self,
@@ -60,6 +61,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
safety_checker: StableDiffusionSafetyChecker,
feature_extractor: CLIPFeatureExtractor,
requires_safety_checker: bool = True,
):
super().__init__()
if hasattr(scheduler.config, "steps_offset") and scheduler.config.steps_offset != 1:
@@ -85,6 +87,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
safety_checker=safety_checker,
feature_extractor=feature_extractor,
)
self.register_to_config(requires_safety_checker=requires_safety_checker)
@property
def components(self) -> Dict[str, Any]:
@@ -121,7 +124,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
def inpaint(
self,
prompt: Union[str, List[str]],
init_image: Union[torch.FloatTensor, PIL.Image.Image],
image: Union[torch.FloatTensor, PIL.Image.Image],
mask_image: Union[torch.FloatTensor, PIL.Image.Image],
strength: float = 0.8,
num_inference_steps: Optional[int] = 50,
@@ -138,7 +141,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
# For more information on how this function works, please see: https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion#diffusers.StableDiffusionImg2ImgPipeline
return StableDiffusionInpaintPipelineLegacy(**self.components)(
prompt=prompt,
init_image=init_image,
image=image,
mask_image=mask_image,
strength=strength,
num_inference_steps=num_inference_steps,
@@ -156,7 +159,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
def img2img(
self,
prompt: Union[str, List[str]],
init_image: Union[torch.FloatTensor, PIL.Image.Image],
image: Union[torch.FloatTensor, PIL.Image.Image],
strength: float = 0.8,
num_inference_steps: Optional[int] = 50,
guidance_scale: Optional[float] = 7.5,
@@ -173,7 +176,7 @@ class StableDiffusionMegaPipeline(DiffusionPipeline):
# For more information on how this function works, please see: https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion#diffusers.StableDiffusionImg2ImgPipeline
return StableDiffusionImg2ImgPipeline(**self.components)(
prompt=prompt,
init_image=init_image,
image=image,
strength=strength,
num_inference_steps=num_inference_steps,
guidance_scale=guidance_scale,

View File

@@ -183,24 +183,6 @@ class TextInpainting(DiffusionPipeline):
return torch.device(module._hf_hook.execution_device)
return self.device
def enable_xformers_memory_efficient_attention(self):
r"""
Enable memory efficient attention as implemented in xformers.
When this option is enabled, you should observe lower GPU memory usage and a potential speed up at inference
time. Speed up at training time is not guaranteed.
Warning: When Memory Efficient Attention and Sliced attention are both enabled, the Memory Efficient Attention
is used.
"""
self.unet.set_use_memory_efficient_attention_xformers(True)
def disable_xformers_memory_efficient_attention(self):
r"""
Disable memory efficient attention as implemented in xformers.
"""
self.unet.set_use_memory_efficient_attention_xformers(False)
@torch.no_grad()
def __call__(
self,

View File

@@ -68,7 +68,7 @@ class WildcardStableDiffusionPipeline(DiffusionPipeline):
Example Usage:
pipe = WildcardStableDiffusionPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
revision="fp16",
torch_dtype=torch.float16,
)
prompt = "__animal__ sitting on a __object__ wearing a __clothing__"

View File

@@ -9,8 +9,18 @@ The `train_dreambooth.py` script shows how to implement the training procedure a
Before running the scripts, make sure to install the library's training dependencies:
**Important**
To make sure you can successfully run the latest versions of the example scripts, we highly recommend **installing from source** and keeping the install up to date as we update the example scripts frequently and install some example-specific requirements. To do this, execute the following steps in a new virtual environment:
```bash
pip install -U -r requirements.txt
git clone https://github.com/huggingface/diffusers
cd diffusers
pip install -e .
```
Then cd in the example folder and run
```bash
pip install -r requirements.txt
```
And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) environment with:
@@ -19,21 +29,20 @@ And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) e
accelerate config
```
### Dog toy example
You need to accept the model license before downloading or using the weights. In this example we'll use model version `v1-4`, so you'll need to visit [its card](https://huggingface.co/CompVis/stable-diffusion-v1-4), read the license and tick the checkbox if you agree.
You have to be a registered user in 🤗 Hugging Face Hub, and you'll also need to use an access token for the code to work. For more information on access tokens, please refer to [this section of the documentation](https://huggingface.co/docs/hub/security-tokens).
Run the following command to authenticate your token
Or for a default accelerate configuration without answering questions about your environment
```bash
huggingface-cli login
accelerate config default
```
If you have already cloned the repo, then you won't need to go through these steps.
Or if your environment doesn't support an interactive shell e.g. a notebook
<br>
```python
from accelerate.utils import write_basic_config
write_basic_config()
```
### Dog toy example
Now let's get our dataset. Download images from [here](https://drive.google.com/drive/folders/1BO_dyz-p65qhBRRMRA4TbZ8qW4rB99JZ) and save them in a directory. This will be our training data.
@@ -63,7 +72,7 @@ accelerate launch train_dreambooth.py \
### Training with prior-preservation loss
Prior-preservation is used to avoid overfitting and language-drift. Refer to the paper to learn more about it. For prior-preservation we first generate images using the model with a class prompt and then use those during training along with our data.
According to the paper, it's recommended to generate `num_epochs * num_samples` images for prior-preservation. 200-300 works well for most cases.
According to the paper, it's recommended to generate `num_epochs * num_samples` images for prior-preservation. 200-300 works well for most cases. The `num_class_images` flag sets the number of images to generate with the class prompt. You can place existing images in `class_data_dir`, and the training script will generate any additional images so that `num_class_images` are present in `class_data_dir` during training time.
```bash
export MODEL_NAME="CompVis/stable-diffusion-v1-4"
@@ -304,3 +313,97 @@ python train_dreambooth_flax.py \
--num_class_images=200 \
--max_train_steps=800
```
### Training with prior-preservation loss
Prior-preservation is used to avoid overfitting and language-drift. Refer to the paper to learn more about it. For prior-preservation we first generate images using the model with a class prompt and then use those during training along with our data.
According to the paper, it's recommended to generate `num_epochs * num_samples` images for prior-preservation. 200-300 works well for most cases.
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export CLASS_DIR="path-to-class-images"
export OUTPUT_DIR="path-to-save-model"
accelerate launch train_dreambooth_inpaint.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--class_data_dir=$CLASS_DIR \
--output_dir=$OUTPUT_DIR \
--with_prior_preservation --prior_loss_weight=1.0 \
--instance_prompt="a photo of sks dog" \
--class_prompt="a photo of dog" \
--resolution=512 \
--train_batch_size=1 \
--gradient_accumulation_steps=1 \
--learning_rate=5e-6 \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--num_class_images=200 \
--max_train_steps=800
```
### Training with gradient checkpointing and 8-bit optimizer:
With the help of gradient checkpointing and the 8-bit optimizer from bitsandbytes it's possible to run train dreambooth on a 16GB GPU.
To install `bitandbytes` please refer to this [readme](https://github.com/TimDettmers/bitsandbytes#requirements--installation).
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export CLASS_DIR="path-to-class-images"
export OUTPUT_DIR="path-to-save-model"
accelerate launch train_dreambooth_inpaint.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--class_data_dir=$CLASS_DIR \
--output_dir=$OUTPUT_DIR \
--with_prior_preservation --prior_loss_weight=1.0 \
--instance_prompt="a photo of sks dog" \
--class_prompt="a photo of dog" \
--resolution=512 \
--train_batch_size=1 \
--gradient_accumulation_steps=2 --gradient_checkpointing \
--use_8bit_adam \
--learning_rate=5e-6 \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--num_class_images=200 \
--max_train_steps=800
```
### Fine-tune text encoder with the UNet.
The script also allows to fine-tune the `text_encoder` along with the `unet`. It's been observed experimentally that fine-tuning `text_encoder` gives much better results especially on faces.
Pass the `--train_text_encoder` argument to the script to enable training `text_encoder`.
___Note: Training text encoder requires more memory, with this option the training won't fit on 16GB GPU. It needs at least 24GB VRAM.___
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export CLASS_DIR="path-to-class-images"
export OUTPUT_DIR="path-to-save-model"
accelerate launch train_dreambooth_inpaint.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--train_text_encoder \
--instance_data_dir=$INSTANCE_DIR \
--class_data_dir=$CLASS_DIR \
--output_dir=$OUTPUT_DIR \
--with_prior_preservation --prior_loss_weight=1.0 \
--instance_prompt="a photo of sks dog" \
--class_prompt="a photo of dog" \
--resolution=512 \
--train_batch_size=1 \
--use_8bit_adam \
--gradient_checkpointing \
--learning_rate=2e-6 \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--num_class_images=200 \
--max_train_steps=800
```

View File

@@ -1,7 +1,6 @@
diffusers>==0.5.0
accelerate
torchvision
transformers>=4.21.0
transformers>=4.25.1
ftfy
tensorboard
modelcards
modelcards

View File

@@ -1,9 +1,8 @@
diffusers>==0.5.1
transformers>=4.21.0
transformers>=4.25.1
flax
optax
torch
torchvision
ftfy
tensorboard
modelcards
modelcards

View File

@@ -3,6 +3,7 @@ import hashlib
import itertools
import math
import os
import warnings
from pathlib import Path
from typing import Optional
@@ -16,6 +17,8 @@ from accelerate.logging import get_logger
from accelerate.utils import set_seed
from diffusers import AutoencoderKL, DDPMScheduler, DiffusionPipeline, UNet2DConditionModel
from diffusers.optimization import get_scheduler
from diffusers.utils import check_min_version
from diffusers.utils.import_utils import is_xformers_available
from huggingface_hub import HfFolder, Repository, whoami
from PIL import Image
from torchvision import transforms
@@ -23,14 +26,17 @@ from tqdm.auto import tqdm
from transformers import AutoTokenizer, PretrainedConfig
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.10.0.dev0")
logger = get_logger(__name__)
def import_model_class_from_model_name_or_path(pretrained_model_name_or_path: str):
def import_model_class_from_model_name_or_path(pretrained_model_name_or_path: str, revision: str):
text_encoder_config = PretrainedConfig.from_pretrained(
pretrained_model_name_or_path,
subfolder="text_encoder",
revision=args.revision,
revision=revision,
)
model_class = text_encoder_config.architectures[0]
@@ -107,8 +113,8 @@ def parse_args(input_args=None):
type=int,
default=100,
help=(
"Minimal class images for prior preservation loss. If not have enough images, additional images will be"
" sampled with class_prompt."
"Minimal class images for prior preservation loss. If there are not enough images already present in"
" class_data_dir, additional images will be sampled with class_prompt."
),
)
parser.add_argument(
@@ -144,7 +150,25 @@ def parse_args(input_args=None):
default=None,
help="Total number of training steps to perform. If provided, overrides num_train_epochs.",
)
parser.add_argument("--save_steps", type=int, default=500, help="Save checkpoint every X updates steps.")
parser.add_argument(
"--checkpointing_steps",
type=int,
default=500,
help=(
"Save a checkpoint of the training state every X updates. These checkpoints can be used both as final"
" checkpoints in case they are better than the last checkpoint, and are also suitable for resuming"
" training using `--resume_from_checkpoint`."
),
)
parser.add_argument(
"--resume_from_checkpoint",
type=str,
default=None,
help=(
"Whether training should be resumed from a previous checkpoint. Use a path saved by"
' `--checkpointing_steps`, or `"latest"` to automatically select the last available checkpoint.'
),
)
parser.add_argument(
"--gradient_accumulation_steps",
type=int,
@@ -233,10 +257,11 @@ def parse_args(input_args=None):
if args.class_prompt is None:
raise ValueError("You must specify prompt for class images.")
else:
# logger is not available yet
if args.class_data_dir is not None:
logger.warning("You need not use --class_data_dir without --with_prior_preservation.")
warnings.warn("You need not use --class_data_dir without --with_prior_preservation.")
if args.class_prompt is not None:
logger.warning("You need not use --class_prompt without --with_prior_preservation.")
warnings.warn("You need not use --class_prompt without --with_prior_preservation.")
return args
@@ -300,9 +325,10 @@ class DreamBoothDataset(Dataset):
example["instance_images"] = self.image_transforms(instance_image)
example["instance_prompt_ids"] = self.tokenizer(
self.instance_prompt,
padding="do_not_pad",
truncation=True,
padding="max_length",
max_length=self.tokenizer.model_max_length,
return_tensors="pt",
).input_ids
if self.class_data_root:
@@ -312,14 +338,37 @@ class DreamBoothDataset(Dataset):
example["class_images"] = self.image_transforms(class_image)
example["class_prompt_ids"] = self.tokenizer(
self.class_prompt,
padding="do_not_pad",
truncation=True,
padding="max_length",
max_length=self.tokenizer.model_max_length,
return_tensors="pt",
).input_ids
return example
def collate_fn(examples, with_prior_preservation=False):
input_ids = [example["instance_prompt_ids"] for example in examples]
pixel_values = [example["instance_images"] for example in examples]
# Concat class and instance examples for prior preservation.
# We do this to avoid doing two forward passes.
if with_prior_preservation:
input_ids += [example["class_prompt_ids"] for example in examples]
pixel_values += [example["class_images"] for example in examples]
pixel_values = torch.stack(pixel_values)
pixel_values = pixel_values.to(memory_format=torch.contiguous_format).float()
input_ids = torch.cat(input_ids, dim=0)
batch = {
"input_ids": input_ids,
"pixel_values": pixel_values,
}
return batch
class PromptDataset(Dataset):
"A simple dataset to prepare the prompts to generate class images on multiple GPUs."
@@ -441,7 +490,7 @@ def main(args):
)
# import correct text encoder class
text_encoder_cls = import_model_class_from_model_name_or_path(args.pretrained_model_name_or_path)
text_encoder_cls = import_model_class_from_model_name_or_path(args.pretrained_model_name_or_path, args.revision)
# Load models and create wrapper for stable diffusion
text_encoder = text_encoder_cls.from_pretrained(
@@ -460,6 +509,15 @@ def main(args):
revision=args.revision,
)
if is_xformers_available():
try:
unet.enable_xformers_memory_efficient_attention()
except Exception as e:
logger.warning(
"Could not enable memory efficient attention. Make sure xformers is installed"
f" correctly and a GPU is available: {e}"
)
vae.requires_grad_(False)
if not args.train_text_encoder:
text_encoder.requires_grad_(False)
@@ -498,7 +556,7 @@ def main(args):
eps=args.adam_epsilon,
)
noise_scheduler = DDPMScheduler.from_config(args.pretrained_model_name_or_path, subfolder="scheduler")
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
train_dataset = DreamBoothDataset(
instance_data_root=args.instance_data_dir,
@@ -510,34 +568,12 @@ def main(args):
center_crop=args.center_crop,
)
def collate_fn(examples):
input_ids = [example["instance_prompt_ids"] for example in examples]
pixel_values = [example["instance_images"] for example in examples]
# Concat class and instance examples for prior preservation.
# We do this to avoid doing two forward passes.
if args.with_prior_preservation:
input_ids += [example["class_prompt_ids"] for example in examples]
pixel_values += [example["class_images"] for example in examples]
pixel_values = torch.stack(pixel_values)
pixel_values = pixel_values.to(memory_format=torch.contiguous_format).float()
input_ids = tokenizer.pad(
{"input_ids": input_ids},
padding="max_length",
max_length=tokenizer.model_max_length,
return_tensors="pt",
).input_ids
batch = {
"input_ids": input_ids,
"pixel_values": pixel_values,
}
return batch
train_dataloader = torch.utils.data.DataLoader(
train_dataset, batch_size=args.train_batch_size, shuffle=True, collate_fn=collate_fn, num_workers=1
train_dataset,
batch_size=args.train_batch_size,
shuffle=True,
collate_fn=lambda examples: collate_fn(examples, args.with_prior_preservation),
num_workers=1,
)
# Scheduler and math around the number of training steps.
@@ -562,6 +598,7 @@ def main(args):
unet, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
unet, optimizer, train_dataloader, lr_scheduler
)
accelerator.register_for_checkpointing(lr_scheduler)
weight_dtype = torch.float32
if accelerator.mixed_precision == "fp16":
@@ -599,16 +636,41 @@ def main(args):
logger.info(f" Total train batch size (w. parallel, distributed & accumulation) = {total_batch_size}")
logger.info(f" Gradient Accumulation steps = {args.gradient_accumulation_steps}")
logger.info(f" Total optimization steps = {args.max_train_steps}")
# Only show the progress bar once on each machine.
progress_bar = tqdm(range(args.max_train_steps), disable=not accelerator.is_local_main_process)
progress_bar.set_description("Steps")
global_step = 0
first_epoch = 0
for epoch in range(args.num_train_epochs):
if args.resume_from_checkpoint:
if args.resume_from_checkpoint != "latest":
path = os.path.basename(args.resume_from_checkpoint)
else:
# Get the mos recent checkpoint
dirs = os.listdir(args.output_dir)
dirs = [d for d in dirs if d.startswith("checkpoint")]
dirs = sorted(dirs, key=lambda x: int(x.split("-")[1]))
path = dirs[-1]
accelerator.print(f"Resuming from checkpoint {path}")
accelerator.load_state(os.path.join(args.output_dir, path))
global_step = int(path.split("-")[1])
resume_global_step = global_step * args.gradient_accumulation_steps
first_epoch = resume_global_step // num_update_steps_per_epoch
resume_step = resume_global_step % num_update_steps_per_epoch
# Only show the progress bar once on each machine.
progress_bar = tqdm(range(global_step, args.max_train_steps), disable=not accelerator.is_local_main_process)
progress_bar.set_description("Steps")
for epoch in range(first_epoch, args.num_train_epochs):
unet.train()
if args.train_text_encoder:
text_encoder.train()
for step, batch in enumerate(train_dataloader):
# Skip steps until we reach the resumed step
if args.resume_from_checkpoint and epoch == first_epoch and step < resume_step:
if step % args.gradient_accumulation_steps == 0:
progress_bar.update(1)
continue
with accelerator.accumulate(unet):
# Convert images to latent space
latents = vae.encode(batch["pixel_values"].to(dtype=weight_dtype)).latent_dist.sample()
@@ -672,16 +734,11 @@ def main(args):
progress_bar.update(1)
global_step += 1
if global_step % args.save_steps == 0:
if global_step % args.checkpointing_steps == 0:
if accelerator.is_main_process:
pipeline = DiffusionPipeline.from_pretrained(
args.pretrained_model_name_or_path,
unet=accelerator.unwrap_model(unet),
text_encoder=accelerator.unwrap_model(text_encoder),
revision=args.revision,
)
save_path = os.path.join(args.output_dir, f"checkpoint-{global_step}")
pipeline.save_pretrained(save_path)
accelerator.save_state(save_path)
logger.info(f"Saved state to {save_path}")
logs = {"loss": loss.detach().item(), "lr": lr_scheduler.get_last_lr()[0]}
progress_bar.set_postfix(**logs)

View File

@@ -23,6 +23,7 @@ from diffusers import (
FlaxUNet2DConditionModel,
)
from diffusers.pipelines.stable_diffusion import FlaxStableDiffusionSafetyChecker
from diffusers.utils import check_min_version
from flax import jax_utils
from flax.training import train_state
from flax.training.common_utils import shard
@@ -33,6 +34,9 @@ from tqdm.auto import tqdm
from transformers import CLIPFeatureExtractor, CLIPTokenizer, FlaxCLIPTextModel, set_seed
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.10.0.dev0")
logger = logging.getLogger(__name__)
@@ -89,8 +93,8 @@ def parse_args():
type=int,
default=100,
help=(
"Minimal class images for prior preservation loss. If not have enough images, additional images will be"
" sampled with class_prompt."
"Minimal class images for prior preservation loss. If there are not enough images already present in"
" class_data_dir, additional images will be sampled with class_prompt."
),
)
parser.add_argument(

View File

@@ -0,0 +1,14 @@
# Research projects
This folder contains various research projects using 🧨 Diffusers.
They are not really maintained by the core maintainers of this library and often require a specific version of Diffusers that is indicated in the requirements file of each folder.
Updating them to the most recent version of the library will require some work.
To use any of them, just run the command
```
pip install -r requirements.txt
```
inside the folder of your choice.
If you need help with any of those, please open an issue where you directly ping the author(s), as indicated at the top of the README of each folder.

View File

@@ -0,0 +1,26 @@
# Dreambooth for the inpainting model
This script was added by @thedarkzeno .
Please note that this script is not actively maintained, you can open an issue and tag @thedarkzeno or @patil-suraj though.
```bash
export MODEL_NAME="runwayml/stable-diffusion-inpainting"
export INSTANCE_DIR="path-to-instance-images"
export OUTPUT_DIR="path-to-save-model"
accelerate launch train_dreambooth_inpaint.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--output_dir=$OUTPUT_DIR \
--instance_prompt="a photo of sks dog" \
--resolution=512 \
--train_batch_size=1 \
--gradient_accumulation_steps=1 \
--learning_rate=5e-6 \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--max_train_steps=400
```
The script is also compatible with prior preservation loss and gradient checkpointing

View File

@@ -0,0 +1,7 @@
diffusers==0.9.0
accelerate
torchvision
transformers>=4.21.0
ftfy
tensorboard
modelcards

View File

@@ -0,0 +1,798 @@
import argparse
import hashlib
import itertools
import math
import os
import random
from pathlib import Path
from typing import Optional
import numpy as np
import torch
import torch.nn.functional as F
import torch.utils.checkpoint
from torch.utils.data import Dataset
from accelerate import Accelerator
from accelerate.logging import get_logger
from accelerate.utils import set_seed
from diffusers import (
AutoencoderKL,
DDPMScheduler,
StableDiffusionInpaintPipeline,
StableDiffusionPipeline,
UNet2DConditionModel,
)
from diffusers.optimization import get_scheduler
from diffusers.utils import check_min_version
from huggingface_hub import HfFolder, Repository, whoami
from PIL import Image, ImageDraw
from torchvision import transforms
from tqdm.auto import tqdm
from transformers import CLIPTextModel, CLIPTokenizer
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.10.0.dev0")
logger = get_logger(__name__)
def prepare_mask_and_masked_image(image, mask):
image = np.array(image.convert("RGB"))
image = image[None].transpose(0, 3, 1, 2)
image = torch.from_numpy(image).to(dtype=torch.float32) / 127.5 - 1.0
mask = np.array(mask.convert("L"))
mask = mask.astype(np.float32) / 255.0
mask = mask[None, None]
mask[mask < 0.5] = 0
mask[mask >= 0.5] = 1
mask = torch.from_numpy(mask)
masked_image = image * (mask < 0.5)
return mask, masked_image
# generate random masks
def random_mask(im_shape, ratio=1, mask_full_image=False):
mask = Image.new("L", im_shape, 0)
draw = ImageDraw.Draw(mask)
size = (random.randint(0, int(im_shape[0] * ratio)), random.randint(0, int(im_shape[1] * ratio)))
# use this to always mask the whole image
if mask_full_image:
size = (int(im_shape[0] * ratio), int(im_shape[1] * ratio))
limits = (im_shape[0] - size[0] // 2, im_shape[1] - size[1] // 2)
center = (random.randint(size[0] // 2, limits[0]), random.randint(size[1] // 2, limits[1]))
draw_type = random.randint(0, 1)
if draw_type == 0 or mask_full_image:
draw.rectangle(
(center[0] - size[0] // 2, center[1] - size[1] // 2, center[0] + size[0] // 2, center[1] + size[1] // 2),
fill=255,
)
else:
draw.ellipse(
(center[0] - size[0] // 2, center[1] - size[1] // 2, center[0] + size[0] // 2, center[1] + size[1] // 2),
fill=255,
)
return mask
def parse_args():
parser = argparse.ArgumentParser(description="Simple example of a training script.")
parser.add_argument(
"--pretrained_model_name_or_path",
type=str,
default=None,
required=True,
help="Path to pretrained model or model identifier from huggingface.co/models.",
)
parser.add_argument(
"--tokenizer_name",
type=str,
default=None,
help="Pretrained tokenizer name or path if not the same as model_name",
)
parser.add_argument(
"--instance_data_dir",
type=str,
default=None,
required=True,
help="A folder containing the training data of instance images.",
)
parser.add_argument(
"--class_data_dir",
type=str,
default=None,
required=False,
help="A folder containing the training data of class images.",
)
parser.add_argument(
"--instance_prompt",
type=str,
default=None,
help="The prompt with identifier specifying the instance",
)
parser.add_argument(
"--class_prompt",
type=str,
default=None,
help="The prompt to specify images in the same class as provided instance images.",
)
parser.add_argument(
"--with_prior_preservation",
default=False,
action="store_true",
help="Flag to add prior preservation loss.",
)
parser.add_argument("--prior_loss_weight", type=float, default=1.0, help="The weight of prior preservation loss.")
parser.add_argument(
"--num_class_images",
type=int,
default=100,
help=(
"Minimal class images for prior preservation loss. If not have enough images, additional images will be"
" sampled with class_prompt."
),
)
parser.add_argument(
"--output_dir",
type=str,
default="text-inversion-model",
help="The output directory where the model predictions and checkpoints will be written.",
)
parser.add_argument("--seed", type=int, default=None, help="A seed for reproducible training.")
parser.add_argument(
"--resolution",
type=int,
default=512,
help=(
"The resolution for input images, all the images in the train/validation dataset will be resized to this"
" resolution"
),
)
parser.add_argument(
"--center_crop", action="store_true", help="Whether to center crop images before resizing to resolution"
)
parser.add_argument("--train_text_encoder", action="store_true", help="Whether to train the text encoder")
parser.add_argument(
"--train_batch_size", type=int, default=4, help="Batch size (per device) for the training dataloader."
)
parser.add_argument(
"--sample_batch_size", type=int, default=4, help="Batch size (per device) for sampling images."
)
parser.add_argument("--num_train_epochs", type=int, default=1)
parser.add_argument(
"--max_train_steps",
type=int,
default=None,
help="Total number of training steps to perform. If provided, overrides num_train_epochs.",
)
parser.add_argument(
"--gradient_accumulation_steps",
type=int,
default=1,
help="Number of updates steps to accumulate before performing a backward/update pass.",
)
parser.add_argument(
"--gradient_checkpointing",
action="store_true",
help="Whether or not to use gradient checkpointing to save memory at the expense of slower backward pass.",
)
parser.add_argument(
"--learning_rate",
type=float,
default=5e-6,
help="Initial learning rate (after the potential warmup period) to use.",
)
parser.add_argument(
"--scale_lr",
action="store_true",
default=False,
help="Scale the learning rate by the number of GPUs, gradient accumulation steps, and batch size.",
)
parser.add_argument(
"--lr_scheduler",
type=str,
default="constant",
help=(
'The scheduler type to use. Choose between ["linear", "cosine", "cosine_with_restarts", "polynomial",'
' "constant", "constant_with_warmup"]'
),
)
parser.add_argument(
"--lr_warmup_steps", type=int, default=500, help="Number of steps for the warmup in the lr scheduler."
)
parser.add_argument(
"--use_8bit_adam", action="store_true", help="Whether or not to use 8-bit Adam from bitsandbytes."
)
parser.add_argument("--adam_beta1", type=float, default=0.9, help="The beta1 parameter for the Adam optimizer.")
parser.add_argument("--adam_beta2", type=float, default=0.999, help="The beta2 parameter for the Adam optimizer.")
parser.add_argument("--adam_weight_decay", type=float, default=1e-2, help="Weight decay to use.")
parser.add_argument("--adam_epsilon", type=float, default=1e-08, help="Epsilon value for the Adam optimizer")
parser.add_argument("--max_grad_norm", default=1.0, type=float, help="Max gradient norm.")
parser.add_argument("--push_to_hub", action="store_true", help="Whether or not to push the model to the Hub.")
parser.add_argument("--hub_token", type=str, default=None, help="The token to use to push to the Model Hub.")
parser.add_argument(
"--hub_model_id",
type=str,
default=None,
help="The name of the repository to keep in sync with the local `output_dir`.",
)
parser.add_argument(
"--logging_dir",
type=str,
default="logs",
help=(
"[TensorBoard](https://www.tensorflow.org/tensorboard) log directory. Will default to"
" *output_dir/runs/**CURRENT_DATETIME_HOSTNAME***."
),
)
parser.add_argument(
"--mixed_precision",
type=str,
default="no",
choices=["no", "fp16", "bf16"],
help=(
"Whether to use mixed precision. Choose"
"between fp16 and bf16 (bfloat16). Bf16 requires PyTorch >= 1.10."
"and an Nvidia Ampere GPU."
),
)
parser.add_argument("--local_rank", type=int, default=-1, help="For distributed training: local_rank")
parser.add_argument(
"--checkpointing_steps",
type=int,
default=500,
help=(
"Save a checkpoint of the training state every X updates. These checkpoints can be used both as final"
" checkpoints in case they are better than the last checkpoint and are suitable for resuming training"
" using `--resume_from_checkpoint`."
),
)
parser.add_argument(
"--resume_from_checkpoint",
type=str,
default=None,
help=(
"Whether training should be resumed from a previous checkpoint. Use a path saved by"
' `--checkpointing_steps`, or `"latest"` to automatically select the last available checkpoint.'
),
)
args = parser.parse_args()
env_local_rank = int(os.environ.get("LOCAL_RANK", -1))
if env_local_rank != -1 and env_local_rank != args.local_rank:
args.local_rank = env_local_rank
if args.instance_data_dir is None:
raise ValueError("You must specify a train data directory.")
if args.with_prior_preservation:
if args.class_data_dir is None:
raise ValueError("You must specify a data directory for class images.")
if args.class_prompt is None:
raise ValueError("You must specify prompt for class images.")
return args
class DreamBoothDataset(Dataset):
"""
A dataset to prepare the instance and class images with the prompts for fine-tuning the model.
It pre-processes the images and the tokenizes prompts.
"""
def __init__(
self,
instance_data_root,
instance_prompt,
tokenizer,
class_data_root=None,
class_prompt=None,
size=512,
center_crop=False,
):
self.size = size
self.center_crop = center_crop
self.tokenizer = tokenizer
self.instance_data_root = Path(instance_data_root)
if not self.instance_data_root.exists():
raise ValueError("Instance images root doesn't exists.")
self.instance_images_path = list(Path(instance_data_root).iterdir())
self.num_instance_images = len(self.instance_images_path)
self.instance_prompt = instance_prompt
self._length = self.num_instance_images
if class_data_root is not None:
self.class_data_root = Path(class_data_root)
self.class_data_root.mkdir(parents=True, exist_ok=True)
self.class_images_path = list(self.class_data_root.iterdir())
self.num_class_images = len(self.class_images_path)
self._length = max(self.num_class_images, self.num_instance_images)
self.class_prompt = class_prompt
else:
self.class_data_root = None
self.image_transforms_resize_and_crop = transforms.Compose(
[
transforms.Resize(size, interpolation=transforms.InterpolationMode.BILINEAR),
transforms.CenterCrop(size) if center_crop else transforms.RandomCrop(size),
]
)
self.image_transforms = transforms.Compose(
[
transforms.ToTensor(),
transforms.Normalize([0.5], [0.5]),
]
)
def __len__(self):
return self._length
def __getitem__(self, index):
example = {}
instance_image = Image.open(self.instance_images_path[index % self.num_instance_images])
if not instance_image.mode == "RGB":
instance_image = instance_image.convert("RGB")
instance_image = self.image_transforms_resize_and_crop(instance_image)
example["PIL_images"] = instance_image
example["instance_images"] = self.image_transforms(instance_image)
example["instance_prompt_ids"] = self.tokenizer(
self.instance_prompt,
padding="do_not_pad",
truncation=True,
max_length=self.tokenizer.model_max_length,
).input_ids
if self.class_data_root:
class_image = Image.open(self.class_images_path[index % self.num_class_images])
if not class_image.mode == "RGB":
class_image = class_image.convert("RGB")
class_image = self.image_transforms_resize_and_crop(class_image)
example["class_images"] = self.image_transforms(class_image)
example["class_PIL_images"] = class_image
example["class_prompt_ids"] = self.tokenizer(
self.class_prompt,
padding="do_not_pad",
truncation=True,
max_length=self.tokenizer.model_max_length,
).input_ids
return example
class PromptDataset(Dataset):
"A simple dataset to prepare the prompts to generate class images on multiple GPUs."
def __init__(self, prompt, num_samples):
self.prompt = prompt
self.num_samples = num_samples
def __len__(self):
return self.num_samples
def __getitem__(self, index):
example = {}
example["prompt"] = self.prompt
example["index"] = index
return example
def get_full_repo_name(model_id: str, organization: Optional[str] = None, token: Optional[str] = None):
if token is None:
token = HfFolder.get_token()
if organization is None:
username = whoami(token)["name"]
return f"{username}/{model_id}"
else:
return f"{organization}/{model_id}"
def main():
args = parse_args()
logging_dir = Path(args.output_dir, args.logging_dir)
accelerator = Accelerator(
gradient_accumulation_steps=args.gradient_accumulation_steps,
mixed_precision=args.mixed_precision,
log_with="tensorboard",
logging_dir=logging_dir,
)
# Currently, it's not possible to do gradient accumulation when training two models with accelerate.accumulate
# This will be enabled soon in accelerate. For now, we don't allow gradient accumulation when training two models.
# TODO (patil-suraj): Remove this check when gradient accumulation with two models is enabled in accelerate.
if args.train_text_encoder and args.gradient_accumulation_steps > 1 and accelerator.num_processes > 1:
raise ValueError(
"Gradient accumulation is not supported when training the text encoder in distributed training. "
"Please set gradient_accumulation_steps to 1. This feature will be supported in the future."
)
if args.seed is not None:
set_seed(args.seed)
if args.with_prior_preservation:
class_images_dir = Path(args.class_data_dir)
if not class_images_dir.exists():
class_images_dir.mkdir(parents=True)
cur_class_images = len(list(class_images_dir.iterdir()))
if cur_class_images < args.num_class_images:
torch_dtype = torch.float16 if accelerator.device.type == "cuda" else torch.float32
pipeline = StableDiffusionInpaintPipeline.from_pretrained(
args.pretrained_model_name_or_path, torch_dtype=torch_dtype, safety_checker=None
)
pipeline.set_progress_bar_config(disable=True)
num_new_images = args.num_class_images - cur_class_images
logger.info(f"Number of class images to sample: {num_new_images}.")
sample_dataset = PromptDataset(args.class_prompt, num_new_images)
sample_dataloader = torch.utils.data.DataLoader(
sample_dataset, batch_size=args.sample_batch_size, num_workers=1
)
sample_dataloader = accelerator.prepare(sample_dataloader)
pipeline.to(accelerator.device)
transform_to_pil = transforms.ToPILImage()
for example in tqdm(
sample_dataloader, desc="Generating class images", disable=not accelerator.is_local_main_process
):
bsz = len(example["prompt"])
fake_images = torch.rand((3, args.resolution, args.resolution))
transform_to_pil = transforms.ToPILImage()
fake_pil_images = transform_to_pil(fake_images)
fake_mask = random_mask((args.resolution, args.resolution), ratio=1, mask_full_image=True)
images = pipeline(prompt=example["prompt"], mask_image=fake_mask, image=fake_pil_images).images
for i, image in enumerate(images):
hash_image = hashlib.sha1(image.tobytes()).hexdigest()
image_filename = class_images_dir / f"{example['index'][i] + cur_class_images}-{hash_image}.jpg"
image.save(image_filename)
del pipeline
if torch.cuda.is_available():
torch.cuda.empty_cache()
# Handle the repository creation
if accelerator.is_main_process:
if args.push_to_hub:
if args.hub_model_id is None:
repo_name = get_full_repo_name(Path(args.output_dir).name, token=args.hub_token)
else:
repo_name = args.hub_model_id
repo = Repository(args.output_dir, clone_from=repo_name)
with open(os.path.join(args.output_dir, ".gitignore"), "w+") as gitignore:
if "step_*" not in gitignore:
gitignore.write("step_*\n")
if "epoch_*" not in gitignore:
gitignore.write("epoch_*\n")
elif args.output_dir is not None:
os.makedirs(args.output_dir, exist_ok=True)
# Load the tokenizer
if args.tokenizer_name:
tokenizer = CLIPTokenizer.from_pretrained(args.tokenizer_name)
elif args.pretrained_model_name_or_path:
tokenizer = CLIPTokenizer.from_pretrained(args.pretrained_model_name_or_path, subfolder="tokenizer")
# Load models and create wrapper for stable diffusion
text_encoder = CLIPTextModel.from_pretrained(args.pretrained_model_name_or_path, subfolder="text_encoder")
vae = AutoencoderKL.from_pretrained(args.pretrained_model_name_or_path, subfolder="vae")
unet = UNet2DConditionModel.from_pretrained(args.pretrained_model_name_or_path, subfolder="unet")
vae.requires_grad_(False)
if not args.train_text_encoder:
text_encoder.requires_grad_(False)
if args.gradient_checkpointing:
unet.enable_gradient_checkpointing()
if args.train_text_encoder:
text_encoder.gradient_checkpointing_enable()
if args.scale_lr:
args.learning_rate = (
args.learning_rate * args.gradient_accumulation_steps * args.train_batch_size * accelerator.num_processes
)
# Use 8-bit Adam for lower memory usage or to fine-tune the model in 16GB GPUs
if args.use_8bit_adam:
try:
import bitsandbytes as bnb
except ImportError:
raise ImportError(
"To use 8-bit Adam, please install the bitsandbytes library: `pip install bitsandbytes`."
)
optimizer_class = bnb.optim.AdamW8bit
else:
optimizer_class = torch.optim.AdamW
params_to_optimize = (
itertools.chain(unet.parameters(), text_encoder.parameters()) if args.train_text_encoder else unet.parameters()
)
optimizer = optimizer_class(
params_to_optimize,
lr=args.learning_rate,
betas=(args.adam_beta1, args.adam_beta2),
weight_decay=args.adam_weight_decay,
eps=args.adam_epsilon,
)
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
train_dataset = DreamBoothDataset(
instance_data_root=args.instance_data_dir,
instance_prompt=args.instance_prompt,
class_data_root=args.class_data_dir if args.with_prior_preservation else None,
class_prompt=args.class_prompt,
tokenizer=tokenizer,
size=args.resolution,
center_crop=args.center_crop,
)
def collate_fn(examples):
input_ids = [example["instance_prompt_ids"] for example in examples]
pixel_values = [example["instance_images"] for example in examples]
# Concat class and instance examples for prior preservation.
# We do this to avoid doing two forward passes.
if args.with_prior_preservation:
input_ids += [example["class_prompt_ids"] for example in examples]
pixel_values += [example["class_images"] for example in examples]
pior_pil = [example["class_PIL_images"] for example in examples]
masks = []
masked_images = []
for example in examples:
pil_image = example["PIL_images"]
# generate a random mask
mask = random_mask(pil_image.size, 1, False)
# prepare mask and masked image
mask, masked_image = prepare_mask_and_masked_image(pil_image, mask)
masks.append(mask)
masked_images.append(masked_image)
if args.with_prior_preservation:
for pil_image in pior_pil:
# generate a random mask
mask = random_mask(pil_image.size, 1, False)
# prepare mask and masked image
mask, masked_image = prepare_mask_and_masked_image(pil_image, mask)
masks.append(mask)
masked_images.append(masked_image)
pixel_values = torch.stack(pixel_values)
pixel_values = pixel_values.to(memory_format=torch.contiguous_format).float()
input_ids = tokenizer.pad({"input_ids": input_ids}, padding=True, return_tensors="pt").input_ids
masks = torch.stack(masks)
masked_images = torch.stack(masked_images)
batch = {"input_ids": input_ids, "pixel_values": pixel_values, "masks": masks, "masked_images": masked_images}
return batch
train_dataloader = torch.utils.data.DataLoader(
train_dataset, batch_size=args.train_batch_size, shuffle=True, collate_fn=collate_fn
)
# Scheduler and math around the number of training steps.
overrode_max_train_steps = False
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
if args.max_train_steps is None:
args.max_train_steps = args.num_train_epochs * num_update_steps_per_epoch
overrode_max_train_steps = True
lr_scheduler = get_scheduler(
args.lr_scheduler,
optimizer=optimizer,
num_warmup_steps=args.lr_warmup_steps * args.gradient_accumulation_steps,
num_training_steps=args.max_train_steps * args.gradient_accumulation_steps,
)
if args.train_text_encoder:
unet, text_encoder, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
unet, text_encoder, optimizer, train_dataloader, lr_scheduler
)
else:
unet, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
unet, optimizer, train_dataloader, lr_scheduler
)
accelerator.register_for_checkpointing(lr_scheduler)
weight_dtype = torch.float32
if args.mixed_precision == "fp16":
weight_dtype = torch.float16
elif args.mixed_precision == "bf16":
weight_dtype = torch.bfloat16
# Move text_encode and vae to gpu.
# For mixed precision training we cast the text_encoder and vae weights to half-precision
# as these models are only used for inference, keeping weights in full precision is not required.
vae.to(accelerator.device, dtype=weight_dtype)
if not args.train_text_encoder:
text_encoder.to(accelerator.device, dtype=weight_dtype)
# We need to recalculate our total training steps as the size of the training dataloader may have changed.
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
if overrode_max_train_steps:
args.max_train_steps = args.num_train_epochs * num_update_steps_per_epoch
# Afterwards we recalculate our number of training epochs
args.num_train_epochs = math.ceil(args.max_train_steps / num_update_steps_per_epoch)
# We need to initialize the trackers we use, and also store our configuration.
# The trackers initializes automatically on the main process.
if accelerator.is_main_process:
accelerator.init_trackers("dreambooth", config=vars(args))
# Train!
total_batch_size = args.train_batch_size * accelerator.num_processes * args.gradient_accumulation_steps
logger.info("***** Running training *****")
logger.info(f" Num examples = {len(train_dataset)}")
logger.info(f" Num batches each epoch = {len(train_dataloader)}")
logger.info(f" Num Epochs = {args.num_train_epochs}")
logger.info(f" Instantaneous batch size per device = {args.train_batch_size}")
logger.info(f" Total train batch size (w. parallel, distributed & accumulation) = {total_batch_size}")
logger.info(f" Gradient Accumulation steps = {args.gradient_accumulation_steps}")
logger.info(f" Total optimization steps = {args.max_train_steps}")
global_step = 0
first_epoch = 0
if args.resume_from_checkpoint:
if args.resume_from_checkpoint != "latest":
path = os.path.basename(args.resume_from_checkpoint)
else:
# Get the most recent checkpoint
dirs = os.listdir(args.output_dir)
dirs = [d for d in dirs if d.startswith("checkpoint")]
dirs = sorted(dirs, key=lambda x: int(x.split("-")[1]))
path = dirs[-1]
accelerator.print(f"Resuming from checkpoint {path}")
accelerator.load_state(os.path.join(args.output_dir, path))
global_step = int(path.split("-")[1])
resume_global_step = global_step * args.gradient_accumulation_steps
first_epoch = resume_global_step // num_update_steps_per_epoch
resume_step = resume_global_step % num_update_steps_per_epoch
# Only show the progress bar once on each machine.
progress_bar = tqdm(range(global_step, args.max_train_steps), disable=not accelerator.is_local_main_process)
progress_bar.set_description("Steps")
for epoch in range(first_epoch, args.num_epochs):
unet.train()
for step, batch in enumerate(train_dataloader):
# Skip steps until we reach the resumed step
if args.resume_from_checkpoint and epoch == first_epoch and step < resume_step:
if step % args.gradient_accumulation_steps == 0:
progress_bar.update(1)
continue
with accelerator.accumulate(unet):
# Convert images to latent space
latents = vae.encode(batch["pixel_values"].to(dtype=weight_dtype)).latent_dist.sample()
latents = latents * 0.18215
# Convert masked images to latent space
masked_latents = vae.encode(
batch["masked_images"].reshape(batch["pixel_values"].shape).to(dtype=weight_dtype)
).latent_dist.sample()
masked_latents = masked_latents * 0.18215
masks = batch["masks"]
# resize the mask to latents shape as we concatenate the mask to the latents
mask = torch.stack(
[
torch.nn.functional.interpolate(mask, size=(args.resolution // 8, args.resolution // 8))
for mask in masks
]
)
mask = mask.reshape(-1, 1, args.resolution // 8, args.resolution // 8)
# Sample noise that we'll add to the latents
noise = torch.randn_like(latents)
bsz = latents.shape[0]
# Sample a random timestep for each image
timesteps = torch.randint(0, noise_scheduler.config.num_train_timesteps, (bsz,), device=latents.device)
timesteps = timesteps.long()
# Add noise to the latents according to the noise magnitude at each timestep
# (this is the forward diffusion process)
noisy_latents = noise_scheduler.add_noise(latents, noise, timesteps)
# concatenate the noised latents with the mask and the masked latents
latent_model_input = torch.cat([noisy_latents, mask, masked_latents], dim=1)
# Get the text embedding for conditioning
encoder_hidden_states = text_encoder(batch["input_ids"])[0]
# Predict the noise residual
noise_pred = unet(latent_model_input, timesteps, encoder_hidden_states).sample
# Get the target for loss depending on the prediction type
if noise_scheduler.config.prediction_type == "epsilon":
target = noise
elif noise_scheduler.config.prediction_type == "v_prediction":
target = noise_scheduler.get_velocity(latents, noise, timesteps)
else:
raise ValueError(f"Unknown prediction type {noise_scheduler.config.prediction_type}")
if args.with_prior_preservation:
# Chunk the noise and noise_pred into two parts and compute the loss on each part separately.
noise_pred, noise_pred_prior = torch.chunk(noise_pred, 2, dim=0)
target, target_prior = torch.chunk(target, 2, dim=0)
# Compute instance loss
loss = F.mse_loss(noise_pred.float(), target.float(), reduction="none").mean([1, 2, 3]).mean()
# Compute prior loss
prior_loss = F.mse_loss(noise_pred_prior.float(), target_prior.float(), reduction="mean")
# Add the prior loss to the instance loss.
loss = loss + args.prior_loss_weight * prior_loss
else:
loss = F.mse_loss(noise_pred.float(), target.float(), reduction="mean")
accelerator.backward(loss)
if accelerator.sync_gradients:
params_to_clip = (
itertools.chain(unet.parameters(), text_encoder.parameters())
if args.train_text_encoder
else unet.parameters()
)
accelerator.clip_grad_norm_(params_to_clip, args.max_grad_norm)
optimizer.step()
lr_scheduler.step()
optimizer.zero_grad()
# Checks if the accelerator has performed an optimization step behind the scenes
if accelerator.sync_gradients:
progress_bar.update(1)
global_step += 1
if global_step % args.checkpointing_steps == 0:
if accelerator.is_main_process:
save_path = os.path.join(args.output_dir, f"checkpoint-{global_step}")
accelerator.save_state(save_path)
logger.info(f"Saved state to {save_path}")
logs = {"loss": loss.detach().item(), "lr": lr_scheduler.get_last_lr()[0]}
progress_bar.set_postfix(**logs)
accelerator.log(logs, step=global_step)
if global_step >= args.max_train_steps:
break
accelerator.wait_for_everyone()
# Create the pipeline using using the trained modules and save it.
if accelerator.is_main_process:
pipeline = StableDiffusionPipeline.from_pretrained(
args.pretrained_model_name_or_path,
unet=accelerator.unwrap_model(unet),
text_encoder=accelerator.unwrap_model(text_encoder),
)
pipeline.save_pretrained(args.output_dir)
if args.push_to_hub:
repo.push_to_hub(commit_message="End of training", blocking=False, auto_lfs_prune=True)
accelerator.end_training()
if __name__ == "__main__":
main()

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## Diffusers examples with Intel optimizations
**This research project is not actively maintained by the diffusers team. For any questions or comments, please make sure to tag @hshen14 .**
This aims to provide diffusers examples with Intel optimizations such as Bfloat16 for training/fine-tuning acceleration and 8-bit integer (INT8) for inference acceleration on Intel platforms.
## Accelerating the fine-tuning for textual inversion
We accelereate the fine-tuning for textual inversion with Intel Extension for PyTorch. The [examples](textual_inversion) enable both single node and multi-node distributed training with Bfloat16 support on Intel Xeon Scalable Processor.
## Accelerating the inference for Stable Diffusion using Bfloat16
We start the inference acceleration with Bfloat16 using Intel Extension for PyTorch. The [script](inference_bf16.py) is generally designed to support standard Stable Diffusion models with Bfloat16 support.
## Accelerating the inference for Stable Diffusion using INT8
Coming soon ...

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import torch
import intel_extension_for_pytorch as ipex
from diffusers import StableDiffusionPipeline
from PIL import Image
def image_grid(imgs, rows, cols):
assert len(imgs) == rows * cols
w, h = imgs[0].size
grid = Image.new("RGB", size=(cols * w, rows * h))
grid_w, grid_h = grid.size
for i, img in enumerate(imgs):
grid.paste(img, box=(i % cols * w, i // cols * h))
return grid
prompt = ["a lovely <dicoo> in red dress and hat, in the snowly and brightly night, with many brighly buildings"]
batch_size = 8
prompt = prompt * batch_size
device = "cpu"
model_id = "path-to-your-trained-model"
model = StableDiffusionPipeline.from_pretrained(model_id)
model = model.to(device)
# to channels last
model.unet = model.unet.to(memory_format=torch.channels_last)
model.vae = model.vae.to(memory_format=torch.channels_last)
model.text_encoder = model.text_encoder.to(memory_format=torch.channels_last)
model.safety_checker = model.safety_checker.to(memory_format=torch.channels_last)
# optimize with ipex
model.unet = ipex.optimize(model.unet.eval(), dtype=torch.bfloat16, inplace=True)
model.vae = ipex.optimize(model.vae.eval(), dtype=torch.bfloat16, inplace=True)
model.text_encoder = ipex.optimize(model.text_encoder.eval(), dtype=torch.bfloat16, inplace=True)
model.safety_checker = ipex.optimize(model.safety_checker.eval(), dtype=torch.bfloat16, inplace=True)
# compute
seed = 666
generator = torch.Generator(device).manual_seed(seed)
with torch.cpu.amp.autocast(enabled=True, dtype=torch.bfloat16):
images = model(prompt, guidance_scale=7.5, num_inference_steps=50, generator=generator).images
# save image
grid = image_grid(images, rows=2, cols=4)
grid.save(model_id + ".png")

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## Textual Inversion fine-tuning example
[Textual inversion](https://arxiv.org/abs/2208.01618) is a method to personalize text2image models like stable diffusion on your own images using just 3-5 examples.
The `textual_inversion.py` script shows how to implement the training procedure and adapt it for stable diffusion.
## Training with Intel Extension for PyTorch
Intel Extension for PyTorch provides the optimizations for faster training and inference on CPUs. You can leverage the training example "textual_inversion.py". Follow the [instructions](https://github.com/huggingface/diffusers/tree/main/examples/textual_inversion) to get the model and [dataset](https://huggingface.co/sd-concepts-library/dicoo2) before running the script.
The example supports both single node and multi-node distributed training:
### Single node training
```bash
export MODEL_NAME="CompVis/stable-diffusion-v1-4"
export DATA_DIR="path-to-dir-containing-dicoo-images"
python textual_inversion.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--train_data_dir=$DATA_DIR \
--learnable_property="object" \
--placeholder_token="<dicoo>" --initializer_token="toy" \
--seed=7 \
--resolution=512 \
--train_batch_size=1 \
--gradient_accumulation_steps=1 \
--max_train_steps=3000 \
--learning_rate=2.5e-03 --scale_lr \
--output_dir="textual_inversion_dicoo"
```
Note: Bfloat16 is available on Intel Xeon Scalable Processors Cooper Lake or Sapphire Rapids. You may not get performance speedup without Bfloat16 support.
### Multi-node distributed training
Before running the scripts, make sure to install the library's training dependencies successfully:
```bash
python -m pip install oneccl_bind_pt==1.13 -f https://developer.intel.com/ipex-whl-stable-cpu
```
```bash
export MODEL_NAME="CompVis/stable-diffusion-v1-4"
export DATA_DIR="path-to-dir-containing-dicoo-images"
oneccl_bindings_for_pytorch_path=$(python -c "from oneccl_bindings_for_pytorch import cwd; print(cwd)")
source $oneccl_bindings_for_pytorch_path/env/setvars.sh
python -m intel_extension_for_pytorch.cpu.launch --distributed \
--hostfile hostfile --nnodes 2 --nproc_per_node 2 textual_inversion.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--train_data_dir=$DATA_DIR \
--learnable_property="object" \
--placeholder_token="<dicoo>" --initializer_token="toy" \
--seed=7 \
--resolution=512 \
--train_batch_size=1 \
--gradient_accumulation_steps=1 \
--max_train_steps=750 \
--learning_rate=2.5e-03 --scale_lr \
--output_dir="textual_inversion_dicoo"
```
The above is a simple distributed training usage on 2 nodes with 2 processes on each node. Add the right hostname or ip address in the "hostfile" and make sure these 2 nodes are reachable from each other. For more details, please refer to the [user guide](https://github.com/intel/torch-ccl).
### Reference
We publish a [Medium blog](https://medium.com/intel-analytics-software/personalized-stable-diffusion-with-few-shot-fine-tuning-on-a-single-cpu-f01a3316b13) on how to create your own Stable Diffusion model on CPUs using textual inversion. Try it out now, if you have interests.

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accelerate
torchvision
transformers>=4.21.0
ftfy
tensorboard
modelcards
intel_extension_for_pytorch>=1.13

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import argparse
import itertools
import math
import os
import random
from pathlib import Path
from typing import Optional
import numpy as np
import torch
import torch.nn.functional as F
import torch.utils.checkpoint
from torch.utils.data import Dataset
import intel_extension_for_pytorch as ipex
import PIL
from accelerate import Accelerator
from accelerate.logging import get_logger
from accelerate.utils import set_seed
from diffusers import AutoencoderKL, DDPMScheduler, PNDMScheduler, StableDiffusionPipeline, UNet2DConditionModel
from diffusers.optimization import get_scheduler
from diffusers.pipelines.stable_diffusion import StableDiffusionSafetyChecker
from diffusers.utils import check_min_version
from huggingface_hub import HfFolder, Repository, whoami
# TODO: remove and import from diffusers.utils when the new version of diffusers is released
from packaging import version
from PIL import Image
from torchvision import transforms
from tqdm.auto import tqdm
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
if version.parse(version.parse(PIL.__version__).base_version) >= version.parse("9.1.0"):
PIL_INTERPOLATION = {
"linear": PIL.Image.Resampling.BILINEAR,
"bilinear": PIL.Image.Resampling.BILINEAR,
"bicubic": PIL.Image.Resampling.BICUBIC,
"lanczos": PIL.Image.Resampling.LANCZOS,
"nearest": PIL.Image.Resampling.NEAREST,
}
else:
PIL_INTERPOLATION = {
"linear": PIL.Image.LINEAR,
"bilinear": PIL.Image.BILINEAR,
"bicubic": PIL.Image.BICUBIC,
"lanczos": PIL.Image.LANCZOS,
"nearest": PIL.Image.NEAREST,
}
# ------------------------------------------------------------------------------
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.10.0.dev0")
logger = get_logger(__name__)
def save_progress(text_encoder, placeholder_token_id, accelerator, args, save_path):
logger.info("Saving embeddings")
learned_embeds = accelerator.unwrap_model(text_encoder).get_input_embeddings().weight[placeholder_token_id]
learned_embeds_dict = {args.placeholder_token: learned_embeds.detach().cpu()}
torch.save(learned_embeds_dict, save_path)
def parse_args():
parser = argparse.ArgumentParser(description="Simple example of a training script.")
parser.add_argument(
"--save_steps",
type=int,
default=500,
help="Save learned_embeds.bin every X updates steps.",
)
parser.add_argument(
"--only_save_embeds",
action="store_true",
default=False,
help="Save only the embeddings for the new concept.",
)
parser.add_argument(
"--pretrained_model_name_or_path",
type=str,
default=None,
required=True,
help="Path to pretrained model or model identifier from huggingface.co/models.",
)
parser.add_argument(
"--revision",
type=str,
default=None,
required=False,
help="Revision of pretrained model identifier from huggingface.co/models.",
)
parser.add_argument(
"--tokenizer_name",
type=str,
default=None,
help="Pretrained tokenizer name or path if not the same as model_name",
)
parser.add_argument(
"--train_data_dir", type=str, default=None, required=True, help="A folder containing the training data."
)
parser.add_argument(
"--placeholder_token",
type=str,
default=None,
required=True,
help="A token to use as a placeholder for the concept.",
)
parser.add_argument(
"--initializer_token", type=str, default=None, required=True, help="A token to use as initializer word."
)
parser.add_argument("--learnable_property", type=str, default="object", help="Choose between 'object' and 'style'")
parser.add_argument("--repeats", type=int, default=100, help="How many times to repeat the training data.")
parser.add_argument(
"--output_dir",
type=str,
default="text-inversion-model",
help="The output directory where the model predictions and checkpoints will be written.",
)
parser.add_argument("--seed", type=int, default=None, help="A seed for reproducible training.")
parser.add_argument(
"--resolution",
type=int,
default=512,
help=(
"The resolution for input images, all the images in the train/validation dataset will be resized to this"
" resolution"
),
)
parser.add_argument(
"--center_crop", action="store_true", help="Whether to center crop images before resizing to resolution"
)
parser.add_argument(
"--train_batch_size", type=int, default=16, help="Batch size (per device) for the training dataloader."
)
parser.add_argument("--num_train_epochs", type=int, default=100)
parser.add_argument(
"--max_train_steps",
type=int,
default=5000,
help="Total number of training steps to perform. If provided, overrides num_train_epochs.",
)
parser.add_argument(
"--gradient_accumulation_steps",
type=int,
default=1,
help="Number of updates steps to accumulate before performing a backward/update pass.",
)
parser.add_argument(
"--learning_rate",
type=float,
default=1e-4,
help="Initial learning rate (after the potential warmup period) to use.",
)
parser.add_argument(
"--scale_lr",
action="store_true",
default=True,
help="Scale the learning rate by the number of GPUs, gradient accumulation steps, and batch size.",
)
parser.add_argument(
"--lr_scheduler",
type=str,
default="constant",
help=(
'The scheduler type to use. Choose between ["linear", "cosine", "cosine_with_restarts", "polynomial",'
' "constant", "constant_with_warmup"]'
),
)
parser.add_argument(
"--lr_warmup_steps", type=int, default=500, help="Number of steps for the warmup in the lr scheduler."
)
parser.add_argument("--adam_beta1", type=float, default=0.9, help="The beta1 parameter for the Adam optimizer.")
parser.add_argument("--adam_beta2", type=float, default=0.999, help="The beta2 parameter for the Adam optimizer.")
parser.add_argument("--adam_weight_decay", type=float, default=1e-2, help="Weight decay to use.")
parser.add_argument("--adam_epsilon", type=float, default=1e-08, help="Epsilon value for the Adam optimizer")
parser.add_argument("--push_to_hub", action="store_true", help="Whether or not to push the model to the Hub.")
parser.add_argument("--hub_token", type=str, default=None, help="The token to use to push to the Model Hub.")
parser.add_argument(
"--hub_model_id",
type=str,
default=None,
help="The name of the repository to keep in sync with the local `output_dir`.",
)
parser.add_argument(
"--logging_dir",
type=str,
default="logs",
help=(
"[TensorBoard](https://www.tensorflow.org/tensorboard) log directory. Will default to"
" *output_dir/runs/**CURRENT_DATETIME_HOSTNAME***."
),
)
parser.add_argument(
"--mixed_precision",
type=str,
default="no",
choices=["no", "fp16", "bf16"],
help=(
"Whether to use mixed precision. Choose"
"between fp16 and bf16 (bfloat16). Bf16 requires PyTorch >= 1.10."
"and an Nvidia Ampere GPU."
),
)
parser.add_argument("--local_rank", type=int, default=-1, help="For distributed training: local_rank")
args = parser.parse_args()
env_local_rank = int(os.environ.get("LOCAL_RANK", -1))
if env_local_rank != -1 and env_local_rank != args.local_rank:
args.local_rank = env_local_rank
if args.train_data_dir is None:
raise ValueError("You must specify a train data directory.")
return args
imagenet_templates_small = [
"a photo of a {}",
"a rendering of a {}",
"a cropped photo of the {}",
"the photo of a {}",
"a photo of a clean {}",
"a photo of a dirty {}",
"a dark photo of the {}",
"a photo of my {}",
"a photo of the cool {}",
"a close-up photo of a {}",
"a bright photo of the {}",
"a cropped photo of a {}",
"a photo of the {}",
"a good photo of the {}",
"a photo of one {}",
"a close-up photo of the {}",
"a rendition of the {}",
"a photo of the clean {}",
"a rendition of a {}",
"a photo of a nice {}",
"a good photo of a {}",
"a photo of the nice {}",
"a photo of the small {}",
"a photo of the weird {}",
"a photo of the large {}",
"a photo of a cool {}",
"a photo of a small {}",
]
imagenet_style_templates_small = [
"a painting in the style of {}",
"a rendering in the style of {}",
"a cropped painting in the style of {}",
"the painting in the style of {}",
"a clean painting in the style of {}",
"a dirty painting in the style of {}",
"a dark painting in the style of {}",
"a picture in the style of {}",
"a cool painting in the style of {}",
"a close-up painting in the style of {}",
"a bright painting in the style of {}",
"a cropped painting in the style of {}",
"a good painting in the style of {}",
"a close-up painting in the style of {}",
"a rendition in the style of {}",
"a nice painting in the style of {}",
"a small painting in the style of {}",
"a weird painting in the style of {}",
"a large painting in the style of {}",
]
class TextualInversionDataset(Dataset):
def __init__(
self,
data_root,
tokenizer,
learnable_property="object", # [object, style]
size=512,
repeats=100,
interpolation="bicubic",
flip_p=0.5,
set="train",
placeholder_token="*",
center_crop=False,
):
self.data_root = data_root
self.tokenizer = tokenizer
self.learnable_property = learnable_property
self.size = size
self.placeholder_token = placeholder_token
self.center_crop = center_crop
self.flip_p = flip_p
self.image_paths = [os.path.join(self.data_root, file_path) for file_path in os.listdir(self.data_root)]
self.num_images = len(self.image_paths)
self._length = self.num_images
if set == "train":
self._length = self.num_images * repeats
self.interpolation = {
"linear": PIL_INTERPOLATION["linear"],
"bilinear": PIL_INTERPOLATION["bilinear"],
"bicubic": PIL_INTERPOLATION["bicubic"],
"lanczos": PIL_INTERPOLATION["lanczos"],
}[interpolation]
self.templates = imagenet_style_templates_small if learnable_property == "style" else imagenet_templates_small
self.flip_transform = transforms.RandomHorizontalFlip(p=self.flip_p)
def __len__(self):
return self._length
def __getitem__(self, i):
example = {}
image = Image.open(self.image_paths[i % self.num_images])
if not image.mode == "RGB":
image = image.convert("RGB")
placeholder_string = self.placeholder_token
text = random.choice(self.templates).format(placeholder_string)
example["input_ids"] = self.tokenizer(
text,
padding="max_length",
truncation=True,
max_length=self.tokenizer.model_max_length,
return_tensors="pt",
).input_ids[0]
# default to score-sde preprocessing
img = np.array(image).astype(np.uint8)
if self.center_crop:
crop = min(img.shape[0], img.shape[1])
h, w, = (
img.shape[0],
img.shape[1],
)
img = img[(h - crop) // 2 : (h + crop) // 2, (w - crop) // 2 : (w + crop) // 2]
image = Image.fromarray(img)
image = image.resize((self.size, self.size), resample=self.interpolation)
image = self.flip_transform(image)
image = np.array(image).astype(np.uint8)
image = (image / 127.5 - 1.0).astype(np.float32)
example["pixel_values"] = torch.from_numpy(image).permute(2, 0, 1)
return example
def get_full_repo_name(model_id: str, organization: Optional[str] = None, token: Optional[str] = None):
if token is None:
token = HfFolder.get_token()
if organization is None:
username = whoami(token)["name"]
return f"{username}/{model_id}"
else:
return f"{organization}/{model_id}"
def freeze_params(params):
for param in params:
param.requires_grad = False
def main():
args = parse_args()
logging_dir = os.path.join(args.output_dir, args.logging_dir)
accelerator = Accelerator(
gradient_accumulation_steps=args.gradient_accumulation_steps,
mixed_precision=args.mixed_precision,
log_with="tensorboard",
logging_dir=logging_dir,
)
# If passed along, set the training seed now.
if args.seed is not None:
set_seed(args.seed)
# Handle the repository creation
if accelerator.is_main_process:
if args.push_to_hub:
if args.hub_model_id is None:
repo_name = get_full_repo_name(Path(args.output_dir).name, token=args.hub_token)
else:
repo_name = args.hub_model_id
repo = Repository(args.output_dir, clone_from=repo_name)
with open(os.path.join(args.output_dir, ".gitignore"), "w+") as gitignore:
if "step_*" not in gitignore:
gitignore.write("step_*\n")
if "epoch_*" not in gitignore:
gitignore.write("epoch_*\n")
elif args.output_dir is not None:
os.makedirs(args.output_dir, exist_ok=True)
# Load the tokenizer and add the placeholder token as a additional special token
if args.tokenizer_name:
tokenizer = CLIPTokenizer.from_pretrained(args.tokenizer_name)
elif args.pretrained_model_name_or_path:
tokenizer = CLIPTokenizer.from_pretrained(args.pretrained_model_name_or_path, subfolder="tokenizer")
# Add the placeholder token in tokenizer
num_added_tokens = tokenizer.add_tokens(args.placeholder_token)
if num_added_tokens == 0:
raise ValueError(
f"The tokenizer already contains the token {args.placeholder_token}. Please pass a different"
" `placeholder_token` that is not already in the tokenizer."
)
# Convert the initializer_token, placeholder_token to ids
token_ids = tokenizer.encode(args.initializer_token, add_special_tokens=False)
# Check if initializer_token is a single token or a sequence of tokens
if len(token_ids) > 1:
raise ValueError("The initializer token must be a single token.")
initializer_token_id = token_ids[0]
placeholder_token_id = tokenizer.convert_tokens_to_ids(args.placeholder_token)
# Load models and create wrapper for stable diffusion
text_encoder = CLIPTextModel.from_pretrained(
args.pretrained_model_name_or_path,
subfolder="text_encoder",
revision=args.revision,
)
vae = AutoencoderKL.from_pretrained(
args.pretrained_model_name_or_path,
subfolder="vae",
revision=args.revision,
)
unet = UNet2DConditionModel.from_pretrained(
args.pretrained_model_name_or_path,
subfolder="unet",
revision=args.revision,
)
# Resize the token embeddings as we are adding new special tokens to the tokenizer
text_encoder.resize_token_embeddings(len(tokenizer))
# Initialise the newly added placeholder token with the embeddings of the initializer token
token_embeds = text_encoder.get_input_embeddings().weight.data
token_embeds[placeholder_token_id] = token_embeds[initializer_token_id]
# Freeze vae and unet
freeze_params(vae.parameters())
freeze_params(unet.parameters())
# Freeze all parameters except for the token embeddings in text encoder
params_to_freeze = itertools.chain(
text_encoder.text_model.encoder.parameters(),
text_encoder.text_model.final_layer_norm.parameters(),
text_encoder.text_model.embeddings.position_embedding.parameters(),
)
freeze_params(params_to_freeze)
if args.scale_lr:
args.learning_rate = (
args.learning_rate * args.gradient_accumulation_steps * args.train_batch_size * accelerator.num_processes
)
# Initialize the optimizer
optimizer = torch.optim.AdamW(
text_encoder.get_input_embeddings().parameters(), # only optimize the embeddings
lr=args.learning_rate,
betas=(args.adam_beta1, args.adam_beta2),
weight_decay=args.adam_weight_decay,
eps=args.adam_epsilon,
)
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
train_dataset = TextualInversionDataset(
data_root=args.train_data_dir,
tokenizer=tokenizer,
size=args.resolution,
placeholder_token=args.placeholder_token,
repeats=args.repeats,
learnable_property=args.learnable_property,
center_crop=args.center_crop,
set="train",
)
train_dataloader = torch.utils.data.DataLoader(train_dataset, batch_size=args.train_batch_size, shuffle=True)
# Scheduler and math around the number of training steps.
overrode_max_train_steps = False
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
if args.max_train_steps is None:
args.max_train_steps = args.num_train_epochs * num_update_steps_per_epoch
overrode_max_train_steps = True
lr_scheduler = get_scheduler(
args.lr_scheduler,
optimizer=optimizer,
num_warmup_steps=args.lr_warmup_steps * args.gradient_accumulation_steps,
num_training_steps=args.max_train_steps * args.gradient_accumulation_steps,
)
text_encoder, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
text_encoder, optimizer, train_dataloader, lr_scheduler
)
# Move vae and unet to device
vae.to(accelerator.device)
unet.to(accelerator.device)
# Keep vae and unet in eval model as we don't train these
vae.eval()
unet.eval()
unet = ipex.optimize(unet, dtype=torch.bfloat16, inplace=True)
vae = ipex.optimize(vae, dtype=torch.bfloat16, inplace=True)
# We need to recalculate our total training steps as the size of the training dataloader may have changed.
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
if overrode_max_train_steps:
args.max_train_steps = args.num_train_epochs * num_update_steps_per_epoch
# Afterwards we recalculate our number of training epochs
args.num_train_epochs = math.ceil(args.max_train_steps / num_update_steps_per_epoch)
# We need to initialize the trackers we use, and also store our configuration.
# The trackers initializes automatically on the main process.
if accelerator.is_main_process:
accelerator.init_trackers("textual_inversion", config=vars(args))
# Train!
total_batch_size = args.train_batch_size * accelerator.num_processes * args.gradient_accumulation_steps
logger.info("***** Running training *****")
logger.info(f" Num examples = {len(train_dataset)}")
logger.info(f" Num Epochs = {args.num_train_epochs}")
logger.info(f" Instantaneous batch size per device = {args.train_batch_size}")
logger.info(f" Total train batch size (w. parallel, distributed & accumulation) = {total_batch_size}")
logger.info(f" Gradient Accumulation steps = {args.gradient_accumulation_steps}")
logger.info(f" Total optimization steps = {args.max_train_steps}")
# Only show the progress bar once on each machine.
progress_bar = tqdm(range(args.max_train_steps), disable=not accelerator.is_local_main_process)
progress_bar.set_description("Steps")
global_step = 0
text_encoder.train()
text_encoder, optimizer = ipex.optimize(text_encoder, optimizer=optimizer, dtype=torch.bfloat16)
for epoch in range(args.num_train_epochs):
for step, batch in enumerate(train_dataloader):
with torch.cpu.amp.autocast(enabled=True, dtype=torch.bfloat16):
with accelerator.accumulate(text_encoder):
# Convert images to latent space
latents = vae.encode(batch["pixel_values"]).latent_dist.sample().detach()
latents = latents * 0.18215
# Sample noise that we'll add to the latents
noise = torch.randn(latents.shape).to(latents.device)
bsz = latents.shape[0]
# Sample a random timestep for each image
timesteps = torch.randint(
0, noise_scheduler.config.num_train_timesteps, (bsz,), device=latents.device
).long()
# Add noise to the latents according to the noise magnitude at each timestep
# (this is the forward diffusion process)
noisy_latents = noise_scheduler.add_noise(latents, noise, timesteps)
# Get the text embedding for conditioning
encoder_hidden_states = text_encoder(batch["input_ids"])[0]
# Predict the noise residual
model_pred = unet(noisy_latents, timesteps, encoder_hidden_states).sample
# Get the target for loss depending on the prediction type
if noise_scheduler.config.prediction_type == "epsilon":
target = noise
elif noise_scheduler.config.prediction_type == "v_prediction":
target = noise_scheduler.get_velocity(latents, noise, timesteps)
else:
raise ValueError(f"Unknown prediction type {noise_scheduler.config.prediction_type}")
loss = F.mse_loss(model_pred, target, reduction="none").mean([1, 2, 3]).mean()
accelerator.backward(loss)
# Zero out the gradients for all token embeddings except the newly added
# embeddings for the concept, as we only want to optimize the concept embeddings
if accelerator.num_processes > 1:
grads = text_encoder.module.get_input_embeddings().weight.grad
else:
grads = text_encoder.get_input_embeddings().weight.grad
# Get the index for tokens that we want to zero the grads for
index_grads_to_zero = torch.arange(len(tokenizer)) != placeholder_token_id
grads.data[index_grads_to_zero, :] = grads.data[index_grads_to_zero, :].fill_(0)
optimizer.step()
lr_scheduler.step()
optimizer.zero_grad()
# Checks if the accelerator has performed an optimization step behind the scenes
if accelerator.sync_gradients:
progress_bar.update(1)
global_step += 1
if global_step % args.save_steps == 0:
save_path = os.path.join(args.output_dir, f"learned_embeds-steps-{global_step}.bin")
save_progress(text_encoder, placeholder_token_id, accelerator, args, save_path)
logs = {"loss": loss.detach().item(), "lr": lr_scheduler.get_last_lr()[0]}
progress_bar.set_postfix(**logs)
accelerator.log(logs, step=global_step)
if global_step >= args.max_train_steps:
break
accelerator.wait_for_everyone()
# Create the pipeline using using the trained modules and save it.
if accelerator.is_main_process:
if args.push_to_hub and args.only_save_embeds:
logger.warn("Enabling full model saving because --push_to_hub=True was specified.")
save_full_model = True
else:
save_full_model = not args.only_save_embeds
if save_full_model:
pipeline = StableDiffusionPipeline(
text_encoder=accelerator.unwrap_model(text_encoder),
vae=vae,
unet=unet,
tokenizer=tokenizer,
scheduler=PNDMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler"),
safety_checker=StableDiffusionSafetyChecker.from_pretrained("CompVis/stable-diffusion-safety-checker"),
feature_extractor=CLIPFeatureExtractor.from_pretrained("openai/clip-vit-base-patch32"),
)
pipeline.save_pretrained(args.output_dir)
# Save the newly trained embeddings
save_path = os.path.join(args.output_dir, "learned_embeds.bin")
save_progress(text_encoder, placeholder_token_id, accelerator, args, save_path)
if args.push_to_hub:
repo.push_to_hub(commit_message="End of training", blocking=False, auto_lfs_prune=True)
accelerator.end_training()
if __name__ == "__main__":
main()

View File

@@ -1,9 +1,12 @@
# Overview
These examples show how to run (Diffuser)[https://arxiv.org/abs/2205.09991] in Diffusers.
There are four scripts,
1. `run_diffuser_locomotion.py` to sample actions and run them in the environment,
2. and `run_diffuser_gen_trajectories.py` to just sample actions from the pre-trained diffusion model.
These examples show how to run [Diffuser](https://arxiv.org/abs/2205.09991) in Diffusers.
There are two ways to use the script, `run_diffuser_locomotion.py`.
The key option is a change of the variable `n_guide_steps`.
When `n_guide_steps=0`, the trajectories are sampled from the diffusion model, but not fine-tuned to maximize reward in the environment.
By default, `n_guide_steps=2` to match the original implementation.
You will need some RL specific requirements to run the examples:

View File

@@ -1,57 +0,0 @@
import d4rl # noqa
import gym
import tqdm
from diffusers.experimental import ValueGuidedRLPipeline
config = dict(
n_samples=64,
horizon=32,
num_inference_steps=20,
n_guide_steps=0,
scale_grad_by_std=True,
scale=0.1,
eta=0.0,
t_grad_cutoff=2,
device="cpu",
)
if __name__ == "__main__":
env_name = "hopper-medium-v2"
env = gym.make(env_name)
pipeline = ValueGuidedRLPipeline.from_pretrained(
"bglick13/hopper-medium-v2-value-function-hor32",
env=env,
)
env.seed(0)
obs = env.reset()
total_reward = 0
total_score = 0
T = 1000
rollout = [obs.copy()]
try:
for t in tqdm.tqdm(range(T)):
# Call the policy
denorm_actions = pipeline(obs, planning_horizon=32)
# execute action in environment
next_observation, reward, terminal, _ = env.step(denorm_actions)
score = env.get_normalized_score(total_reward)
# update return
total_reward += reward
total_score += score
print(
f"Step: {t}, Reward: {reward}, Total Reward: {total_reward}, Score: {score}, Total Score:"
f" {total_score}"
)
# save observations for rendering
rollout.append(next_observation.copy())
obs = next_observation
except KeyboardInterrupt:
pass
print(f"Total reward: {total_reward}")

View File

@@ -8,7 +8,7 @@ config = dict(
n_samples=64,
horizon=32,
num_inference_steps=20,
n_guide_steps=2,
n_guide_steps=2, # can set to 0 for faster sampling, does not use value network
scale_grad_by_std=True,
scale=0.1,
eta=0.0,
@@ -40,6 +40,7 @@ if __name__ == "__main__":
# execute action in environment
next_observation, reward, terminal, _ = env.step(denorm_actions)
score = env.get_normalized_score(total_reward)
# update return
total_reward += reward
total_score += score
@@ -47,6 +48,7 @@ if __name__ == "__main__":
f"Step: {t}, Reward: {reward}, Total Reward: {total_reward}, Score: {score}, Total Score:"
f" {total_score}"
)
# save observations for rendering
rollout.append(next_observation.copy())

View File

@@ -12,9 +12,18 @@ ___This script is experimental. The script fine-tunes the whole model and often
Before running the scripts, make sure to install the library's training dependencies:
**Important**
To make sure you can successfully run the latest versions of the example scripts, we highly recommend **installing from source** and keeping the install up to date as we update the example scripts frequently and install some example-specific requirements. To do this, execute the following steps in a new virtual environment:
```bash
pip install git+https://github.com/huggingface/diffusers.git
pip install -U -r requirements.txt
git clone https://github.com/huggingface/diffusers
cd diffusers
pip install .
```
Then cd in the example folder and run
```bash
pip install -r requirements.txt
```
And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) environment with:

View File

@@ -1,7 +1,7 @@
diffusers==0.4.1
accelerate
torchvision
transformers>=4.21.0
transformers>=4.25.1
datasets
ftfy
tensorboard
modelcards
modelcards

View File

@@ -1,9 +1,9 @@
diffusers>==0.5.1
transformers>=4.21.0
transformers>=4.25.1
datasets
flax
optax
torch
torchvision
ftfy
tensorboard
modelcards
modelcards

View File

@@ -17,12 +17,17 @@ from accelerate.utils import set_seed
from datasets import load_dataset
from diffusers import AutoencoderKL, DDPMScheduler, StableDiffusionPipeline, UNet2DConditionModel
from diffusers.optimization import get_scheduler
from diffusers.utils import check_min_version
from diffusers.utils.import_utils import is_xformers_available
from huggingface_hub import HfFolder, Repository, whoami
from torchvision import transforms
from tqdm.auto import tqdm
from transformers import CLIPTextModel, CLIPTokenizer
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.10.0.dev0")
logger = get_logger(__name__)
@@ -211,6 +216,24 @@ def parse_args():
),
)
parser.add_argument("--local_rank", type=int, default=-1, help="For distributed training: local_rank")
parser.add_argument(
"--checkpointing_steps",
type=int,
default=500,
help=(
"Save a checkpoint of the training state every X updates. These checkpoints are only suitable for resuming"
" training using `--resume_from_checkpoint`."
),
)
parser.add_argument(
"--resume_from_checkpoint",
type=str,
default=None,
help=(
"Whether training should be resumed from a previous checkpoint. Use a path saved by"
' `--checkpointing_steps`, or `"latest"` to automatically select the last available checkpoint.'
),
)
args = parser.parse_args()
env_local_rank = int(os.environ.get("LOCAL_RANK", -1))
@@ -360,6 +383,15 @@ def main():
revision=args.revision,
)
if is_xformers_available():
try:
unet.enable_xformers_memory_efficient_attention()
except Exception as e:
logger.warning(
"Could not enable memory efficient attention. Make sure xformers is installed"
f" correctly and a GPU is available: {e}"
)
# Freeze vae and text_encoder
vae.requires_grad_(False)
text_encoder.requires_grad_(False)
@@ -392,7 +424,7 @@ def main():
weight_decay=args.adam_weight_decay,
eps=args.adam_epsilon,
)
noise_scheduler = DDPMScheduler.from_config(args.pretrained_model_name_or_path, subfolder="scheduler")
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
# Get the datasets: you can either provide your own training and evaluation files (see below)
# or specify a Dataset from the hub (the dataset will be downloaded automatically from the datasets Hub).
@@ -514,6 +546,7 @@ def main():
unet, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
unet, optimizer, train_dataloader, lr_scheduler
)
accelerator.register_for_checkpointing(lr_scheduler)
weight_dtype = torch.float32
if accelerator.mixed_precision == "fp16":
@@ -553,16 +586,40 @@ def main():
logger.info(f" Total train batch size (w. parallel, distributed & accumulation) = {total_batch_size}")
logger.info(f" Gradient Accumulation steps = {args.gradient_accumulation_steps}")
logger.info(f" Total optimization steps = {args.max_train_steps}")
global_step = 0
first_epoch = 0
if args.resume_from_checkpoint:
if args.resume_from_checkpoint != "latest":
path = os.path.basename(args.resume_from_checkpoint)
else:
# Get the most recent checkpoint
dirs = os.listdir(args.output_dir)
dirs = [d for d in dirs if d.startswith("checkpoint")]
dirs = sorted(dirs, key=lambda x: int(x.split("-")[1]))
path = dirs[-1]
accelerator.print(f"Resuming from checkpoint {path}")
accelerator.load_state(os.path.join(args.output_dir, path))
global_step = int(path.split("-")[1])
resume_global_step = global_step * args.gradient_accumulation_steps
first_epoch = resume_global_step // num_update_steps_per_epoch
resume_step = resume_global_step % num_update_steps_per_epoch
# Only show the progress bar once on each machine.
progress_bar = tqdm(range(args.max_train_steps), disable=not accelerator.is_local_main_process)
progress_bar = tqdm(range(global_step, args.max_train_steps), disable=not accelerator.is_local_main_process)
progress_bar.set_description("Steps")
global_step = 0
for epoch in range(args.num_train_epochs):
for epoch in range(first_epoch, args.num_train_epochs):
unet.train()
train_loss = 0.0
for step, batch in enumerate(train_dataloader):
# Skip steps until we reach the resumed step
if args.resume_from_checkpoint and epoch == first_epoch and step < resume_step:
if step % args.gradient_accumulation_steps == 0:
progress_bar.update(1)
continue
with accelerator.accumulate(unet):
# Convert images to latent space
latents = vae.encode(batch["pixel_values"].to(weight_dtype)).latent_dist.sample()
@@ -615,6 +672,12 @@ def main():
accelerator.log({"train_loss": train_loss}, step=global_step)
train_loss = 0.0
if global_step % args.checkpointing_steps == 0:
if accelerator.is_main_process:
save_path = os.path.join(args.output_dir, f"checkpoint-{global_step}")
accelerator.save_state(save_path)
logger.info(f"Saved state to {save_path}")
logs = {"step_loss": loss.detach().item(), "lr": lr_scheduler.get_last_lr()[0]}
progress_bar.set_postfix(**logs)

View File

@@ -23,6 +23,7 @@ from diffusers import (
FlaxUNet2DConditionModel,
)
from diffusers.pipelines.stable_diffusion import FlaxStableDiffusionSafetyChecker
from diffusers.utils import check_min_version
from flax import jax_utils
from flax.training import train_state
from flax.training.common_utils import shard
@@ -32,6 +33,9 @@ from tqdm.auto import tqdm
from transformers import CLIPFeatureExtractor, CLIPTokenizer, FlaxCLIPTextModel, set_seed
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.10.0.dev0")
logger = logging.getLogger(__name__)

View File

@@ -16,8 +16,18 @@ Colab for inference
Before running the scripts, make sure to install the library's training dependencies:
**Important**
To make sure you can successfully run the latest versions of the example scripts, we highly recommend **installing from source** and keeping the install up to date as we update the example scripts frequently and install some example-specific requirements. To do this, execute the following steps in a new virtual environment:
```bash
pip install diffusers"[training]" accelerate "transformers>=4.21.0"
git clone https://github.com/huggingface/diffusers
cd diffusers
pip install .
```
Then cd in the example folder and run
```bash
pip install -r requirements.txt
```
And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) environment with:

View File

@@ -1,3 +1,6 @@
accelerate
torchvision
transformers>=4.21.0
transformers>=4.25.1
ftfy
tensorboard
modelcards

View File

@@ -1,9 +1,8 @@
diffusers>==0.5.1
transformers>=4.21.0
transformers>=4.25.1
flax
optax
torch
torchvision
ftfy
tensorboard
modelcards
modelcards

View File

@@ -16,8 +16,11 @@ import PIL
from accelerate import Accelerator
from accelerate.logging import get_logger
from accelerate.utils import set_seed
from diffusers import AutoencoderKL, DDPMScheduler, StableDiffusionPipeline, UNet2DConditionModel
from diffusers import AutoencoderKL, DDPMScheduler, PNDMScheduler, StableDiffusionPipeline, UNet2DConditionModel
from diffusers.optimization import get_scheduler
from diffusers.pipelines.stable_diffusion import StableDiffusionSafetyChecker
from diffusers.utils import check_min_version
from diffusers.utils.import_utils import is_xformers_available
from huggingface_hub import HfFolder, Repository, whoami
# TODO: remove and import from diffusers.utils when the new version of diffusers is released
@@ -25,7 +28,7 @@ from packaging import version
from PIL import Image
from torchvision import transforms
from tqdm.auto import tqdm
from transformers import CLIPTextModel, CLIPTokenizer
from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
if version.parse(version.parse(PIL.__version__).base_version) >= version.parse("9.1.0"):
@@ -47,6 +50,10 @@ else:
# ------------------------------------------------------------------------------
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.10.0.dev0")
logger = get_logger(__name__)
@@ -65,6 +72,12 @@ def parse_args():
default=500,
help="Save learned_embeds.bin every X updates steps.",
)
parser.add_argument(
"--only_save_embeds",
action="store_true",
default=False,
help="Save only the embeddings for the new concept.",
)
parser.add_argument(
"--pretrained_model_name_or_path",
type=str,
@@ -192,6 +205,24 @@ def parse_args():
),
)
parser.add_argument("--local_rank", type=int, default=-1, help="For distributed training: local_rank")
parser.add_argument(
"--checkpointing_steps",
type=int,
default=500,
help=(
"Save a checkpoint of the training state every X updates. These checkpoints are only suitable for resuming"
" training using `--resume_from_checkpoint`."
),
)
parser.add_argument(
"--resume_from_checkpoint",
type=str,
default=None,
help=(
"Whether training should be resumed from a previous checkpoint. Use a path saved by"
' `--checkpointing_steps`, or `"latest"` to automatically select the last available checkpoint.'
),
)
args = parser.parse_args()
env_local_rank = int(os.environ.get("LOCAL_RANK", -1))
@@ -427,6 +458,15 @@ def main():
revision=args.revision,
)
if is_xformers_available():
try:
unet.enable_xformers_memory_efficient_attention()
except Exception as e:
logger.warning(
"Could not enable memory efficient attention. Make sure xformers is installed"
f" correctly and a GPU is available: {e}"
)
# Resize the token embeddings as we are adding new special tokens to the tokenizer
text_encoder.resize_token_embeddings(len(tokenizer))
@@ -459,7 +499,7 @@ def main():
eps=args.adam_epsilon,
)
noise_scheduler = DDPMScheduler.from_config(args.pretrained_model_name_or_path, subfolder="scheduler")
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
train_dataset = TextualInversionDataset(
data_root=args.train_data_dir,
@@ -490,10 +530,17 @@ def main():
text_encoder, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
text_encoder, optimizer, train_dataloader, lr_scheduler
)
accelerator.register_for_checkpointing(lr_scheduler)
weight_dtype = torch.float32
if accelerator.mixed_precision == "fp16":
weight_dtype = torch.float16
elif accelerator.mixed_precision == "bf16":
weight_dtype = torch.bfloat16
# Move vae and unet to device
vae.to(accelerator.device)
unet.to(accelerator.device)
unet.to(accelerator.device, dtype=weight_dtype)
vae.to(accelerator.device, dtype=weight_dtype)
# Keep vae and unet in eval model as we don't train these
vae.eval()
@@ -521,21 +568,49 @@ def main():
logger.info(f" Total train batch size (w. parallel, distributed & accumulation) = {total_batch_size}")
logger.info(f" Gradient Accumulation steps = {args.gradient_accumulation_steps}")
logger.info(f" Total optimization steps = {args.max_train_steps}")
# Only show the progress bar once on each machine.
progress_bar = tqdm(range(args.max_train_steps), disable=not accelerator.is_local_main_process)
progress_bar.set_description("Steps")
global_step = 0
first_epoch = 0
for epoch in range(args.num_train_epochs):
if args.resume_from_checkpoint:
if args.resume_from_checkpoint != "latest":
path = os.path.basename(args.resume_from_checkpoint)
else:
# Get the most recent checkpoint
dirs = os.listdir(args.output_dir)
dirs = [d for d in dirs if d.startswith("checkpoint")]
dirs = sorted(dirs, key=lambda x: int(x.split("-")[1]))
path = dirs[-1]
accelerator.print(f"Resuming from checkpoint {path}")
accelerator.load_state(os.path.join(args.output_dir, path))
global_step = int(path.split("-")[1])
resume_global_step = global_step * args.gradient_accumulation_steps
first_epoch = resume_global_step // num_update_steps_per_epoch
resume_step = resume_global_step % num_update_steps_per_epoch
# Only show the progress bar once on each machine.
progress_bar = tqdm(range(global_step, args.max_train_steps), disable=not accelerator.is_local_main_process)
progress_bar.set_description("Steps")
# keep original embeddings as reference
orig_embeds_params = text_encoder.get_input_embeddings().weight.data.clone()
for epoch in range(first_epoch, args.num_train_epochs):
text_encoder.train()
for step, batch in enumerate(train_dataloader):
# Skip steps until we reach the resumed step
if args.resume_from_checkpoint and epoch == first_epoch and step < resume_step:
if step % args.gradient_accumulation_steps == 0:
progress_bar.update(1)
continue
with accelerator.accumulate(text_encoder):
# Convert images to latent space
latents = vae.encode(batch["pixel_values"]).latent_dist.sample().detach()
latents = vae.encode(batch["pixel_values"].to(dtype=weight_dtype)).latent_dist.sample().detach()
latents = latents * 0.18215
# Sample noise that we'll add to the latents
noise = torch.randn(latents.shape).to(latents.device)
noise = torch.randn(latents.shape).to(latents.device).to(dtype=weight_dtype)
bsz = latents.shape[0]
# Sample a random timestep for each image
timesteps = torch.randint(
@@ -547,7 +622,7 @@ def main():
noisy_latents = noise_scheduler.add_noise(latents, noise, timesteps)
# Get the text embedding for conditioning
encoder_hidden_states = text_encoder(batch["input_ids"])[0]
encoder_hidden_states = text_encoder(batch["input_ids"])[0].to(dtype=weight_dtype)
# Predict the noise residual
model_pred = unet(noisy_latents, timesteps, encoder_hidden_states).sample
@@ -560,23 +635,18 @@ def main():
else:
raise ValueError(f"Unknown prediction type {noise_scheduler.config.prediction_type}")
loss = F.mse_loss(model_pred, target, reduction="none").mean([1, 2, 3]).mean()
loss = F.mse_loss(model_pred.float(), target.float(), reduction="none").mean([1, 2, 3]).mean()
accelerator.backward(loss)
# Zero out the gradients for all token embeddings except the newly added
# embeddings for the concept, as we only want to optimize the concept embeddings
if accelerator.num_processes > 1:
grads = text_encoder.module.get_input_embeddings().weight.grad
else:
grads = text_encoder.get_input_embeddings().weight.grad
# Get the index for tokens that we want to zero the grads for
index_grads_to_zero = torch.arange(len(tokenizer)) != placeholder_token_id
grads.data[index_grads_to_zero, :] = grads.data[index_grads_to_zero, :].fill_(0)
optimizer.step()
lr_scheduler.step()
optimizer.zero_grad()
# Let's make sure we don't update any embedding weights besides the newly added token
index_no_updates = torch.arange(len(tokenizer)) != placeholder_token_id
with torch.no_grad():
text_encoder.get_input_embeddings().weight[index_no_updates] = orig_embeds_params[index_no_updates]
# Checks if the accelerator has performed an optimization step behind the scenes
if accelerator.sync_gradients:
progress_bar.update(1)
@@ -585,6 +655,12 @@ def main():
save_path = os.path.join(args.output_dir, f"learned_embeds-steps-{global_step}.bin")
save_progress(text_encoder, placeholder_token_id, accelerator, args, save_path)
if global_step % args.checkpointing_steps == 0:
if accelerator.is_main_process:
save_path = os.path.join(args.output_dir, f"checkpoint-{global_step}")
accelerator.save_state(save_path)
logger.info(f"Saved state to {save_path}")
logs = {"loss": loss.detach().item(), "lr": lr_scheduler.get_last_lr()[0]}
progress_bar.set_postfix(**logs)
accelerator.log(logs, step=global_step)
@@ -596,16 +672,23 @@ def main():
# Create the pipeline using using the trained modules and save it.
if accelerator.is_main_process:
pipeline = StableDiffusionPipeline.from_pretrained(
args.pretrained_model_name_or_path,
text_encoder=accelerator.unwrap_model(text_encoder),
tokenizer=tokenizer,
vae=vae,
unet=unet,
revision=args.revision,
)
pipeline.save_pretrained(args.output_dir)
# Also save the newly trained embeddings
if args.push_to_hub and args.only_save_embeds:
logger.warn("Enabling full model saving because --push_to_hub=True was specified.")
save_full_model = True
else:
save_full_model = not args.only_save_embeds
if save_full_model:
pipeline = StableDiffusionPipeline(
text_encoder=accelerator.unwrap_model(text_encoder),
vae=vae,
unet=unet,
tokenizer=tokenizer,
scheduler=PNDMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler"),
safety_checker=StableDiffusionSafetyChecker.from_pretrained("CompVis/stable-diffusion-safety-checker"),
feature_extractor=CLIPFeatureExtractor.from_pretrained("openai/clip-vit-base-patch32"),
)
pipeline.save_pretrained(args.output_dir)
# Save the newly trained embeddings
save_path = os.path.join(args.output_dir, "learned_embeds.bin")
save_progress(text_encoder, placeholder_token_id, accelerator, args, save_path)

View File

@@ -24,6 +24,7 @@ from diffusers import (
FlaxUNet2DConditionModel,
)
from diffusers.pipelines.stable_diffusion import FlaxStableDiffusionSafetyChecker
from diffusers.utils import check_min_version
from flax import jax_utils
from flax.training import train_state
from flax.training.common_utils import shard
@@ -55,6 +56,9 @@ else:
}
# ------------------------------------------------------------------------------
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.10.0.dev0")
logger = logging.getLogger(__name__)

View File

@@ -6,10 +6,21 @@ Creating a training image set is [described in a different document](https://hug
Before running the scripts, make sure to install the library's training dependencies:
**Important**
To make sure you can successfully run the latest versions of the example scripts, we highly recommend **installing from source** and keeping the install up to date as we update the example scripts frequently and install some example-specific requirements. To do this, execute the following steps in a new virtual environment:
```bash
pip install diffusers[training] accelerate datasets tensorboard
git clone https://github.com/huggingface/diffusers
cd diffusers
pip install .
```
Then cd in the example folder and run
```bash
pip install -r requirements.txt
```
And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) environment with:
```bash

View File

@@ -11,12 +11,11 @@ import torch.nn.functional as F
from accelerate import Accelerator
from accelerate.logging import get_logger
from datasets import load_dataset
from diffusers import DDPMPipeline, DDPMScheduler, UNet2DModel, __version__
from diffusers import DDPMPipeline, DDPMScheduler, UNet2DModel
from diffusers.optimization import get_scheduler
from diffusers.training_utils import EMAModel
from diffusers.utils import deprecate
from diffusers.utils import check_min_version
from huggingface_hub import HfFolder, Repository, whoami
from packaging import version
from torchvision.transforms import (
CenterCrop,
Compose,
@@ -29,8 +28,11 @@ from torchvision.transforms import (
from tqdm.auto import tqdm
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.10.0.dev0")
logger = get_logger(__name__)
diffusers_version = version.parse(version.parse(__version__).base_version)
def _extract_into_tensor(arr, timesteps, broadcast_shape):
@@ -171,6 +173,16 @@ def parse_args():
parser.add_argument(
"--hub_private_repo", action="store_true", help="Whether or not to create a private repository."
)
parser.add_argument(
"--logger",
type=str,
default="tensorboard",
choices=["tensorboard", "wandb"],
help=(
"Whether to use [tensorboard](https://www.tensorflow.org/tensorboard) or [wandb](https://www.wandb.ai)"
" for experiment tracking and logging of model metrics and model checkpoints"
),
)
parser.add_argument(
"--logging_dir",
type=str,
@@ -192,7 +204,6 @@ def parse_args():
"and an Nvidia Ampere GPU."
),
)
parser.add_argument(
"--prediction_type",
type=str,
@@ -200,9 +211,26 @@ def parse_args():
choices=["epsilon", "sample"],
help="Whether the model should predict the 'epsilon'/noise error or directly the reconstructed image 'x0'.",
)
parser.add_argument("--ddpm_num_steps", type=int, default=1000)
parser.add_argument("--ddpm_beta_schedule", type=str, default="linear")
parser.add_argument(
"--checkpointing_steps",
type=int,
default=500,
help=(
"Save a checkpoint of the training state every X updates. These checkpoints are only suitable for resuming"
" training using `--resume_from_checkpoint`."
),
)
parser.add_argument(
"--resume_from_checkpoint",
type=str,
default=None,
help=(
"Whether training should be resumed from a previous checkpoint. Use a path saved by"
' `--checkpointing_steps`, or `"latest"` to automatically select the last available checkpoint.'
),
)
args = parser.parse_args()
env_local_rank = int(os.environ.get("LOCAL_RANK", -1))
@@ -230,7 +258,7 @@ def main(args):
accelerator = Accelerator(
gradient_accumulation_steps=args.gradient_accumulation_steps,
mixed_precision=args.mixed_precision,
log_with="tensorboard",
log_with=args.logger,
logging_dir=logging_dir,
)
@@ -317,6 +345,7 @@ def main(args):
model, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
model, optimizer, train_dataloader, lr_scheduler
)
accelerator.register_for_checkpointing(lr_scheduler)
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
@@ -349,11 +378,36 @@ def main(args):
accelerator.init_trackers(run)
global_step = 0
for epoch in range(args.num_epochs):
first_epoch = 0
if args.resume_from_checkpoint:
if args.resume_from_checkpoint != "latest":
path = os.path.basename(args.resume_from_checkpoint)
else:
# Get the most recent checkpoint
dirs = os.listdir(args.output_dir)
dirs = [d for d in dirs if d.startswith("checkpoint")]
dirs = sorted(dirs, key=lambda x: int(x.split("-")[1]))
path = dirs[-1]
accelerator.print(f"Resuming from checkpoint {path}")
accelerator.load_state(os.path.join(args.output_dir, path))
global_step = int(path.split("-")[1])
resume_global_step = global_step * args.gradient_accumulation_steps
first_epoch = resume_global_step // num_update_steps_per_epoch
resume_step = resume_global_step % num_update_steps_per_epoch
for epoch in range(first_epoch, args.num_epochs):
model.train()
progress_bar = tqdm(total=num_update_steps_per_epoch, disable=not accelerator.is_local_main_process)
progress_bar.set_description(f"Epoch {epoch}")
for step, batch in enumerate(train_dataloader):
# Skip steps until we reach the resumed step
if args.resume_from_checkpoint and epoch == first_epoch and step < resume_step:
if step % args.gradient_accumulation_steps == 0:
progress_bar.update(1)
continue
clean_images = batch["input"]
# Sample noise that we'll add to the images
noise = torch.randn(clean_images.shape).to(clean_images.device)
@@ -400,6 +454,12 @@ def main(args):
progress_bar.update(1)
global_step += 1
if global_step % args.checkpointing_steps == 0:
if accelerator.is_main_process:
save_path = os.path.join(args.output_dir, f"checkpoint-{global_step}")
accelerator.save_state(save_path)
logger.info(f"Saved state to {save_path}")
logs = {"loss": loss.detach().item(), "lr": lr_scheduler.get_last_lr()[0], "step": global_step}
if args.use_ema:
logs["ema_decay"] = ema_model.decay
@@ -417,11 +477,7 @@ def main(args):
scheduler=noise_scheduler,
)
deprecate("todo: remove this check", "0.10.0", "when the most used version is >= 0.8.0")
if diffusers_version < version.parse("0.8.0"):
generator = torch.manual_seed(0)
else:
generator = torch.Generator(device=pipeline.device).manual_seed(0)
generator = torch.Generator(device=pipeline.device).manual_seed(0)
# run pipeline in inference (sample random noise and denoise)
images = pipeline(
generator=generator,
@@ -431,9 +487,11 @@ def main(args):
# denormalize the images and save to tensorboard
images_processed = (images * 255).round().astype("uint8")
accelerator.trackers[0].writer.add_images(
"test_samples", images_processed.transpose(0, 3, 1, 2), epoch
)
if args.logger == "tensorboard":
accelerator.get_tracker("tensorboard").add_images(
"test_samples", images_processed.transpose(0, 3, 1, 2), epoch
)
if epoch % args.save_model_epochs == 0 or epoch == args.num_epochs - 1:
# save the model

View File

@@ -1,6 +1,9 @@
import argparse
import inspect
import math
import os
from pathlib import Path
from typing import Optional
import torch
import torch.nn.functional as F
@@ -9,9 +12,10 @@ from accelerate import Accelerator
from accelerate.logging import get_logger
from datasets import load_dataset
from diffusers import DDPMPipeline, DDPMScheduler, UNet2DModel
from diffusers.hub_utils import init_git_repo, push_to_hub
from diffusers.optimization import get_scheduler
from diffusers.training_utils import EMAModel
from diffusers.utils import check_min_version
from huggingface_hub import HfFolder, Repository, whoami
from onnxruntime.training.ortmodule import ORTModule
from torchvision.transforms import (
CenterCrop,
@@ -25,9 +29,205 @@ from torchvision.transforms import (
from tqdm.auto import tqdm
# Will error if the minimal version of diffusers is not installed. Remove at your own risks.
check_min_version("0.10.0.dev0")
logger = get_logger(__name__)
def _extract_into_tensor(arr, timesteps, broadcast_shape):
"""
Extract values from a 1-D numpy array for a batch of indices.
:param arr: the 1-D numpy array.
:param timesteps: a tensor of indices into the array to extract.
:param broadcast_shape: a larger shape of K dimensions with the batch
dimension equal to the length of timesteps.
:return: a tensor of shape [batch_size, 1, ...] where the shape has K dims.
"""
if not isinstance(arr, torch.Tensor):
arr = torch.from_numpy(arr)
res = arr[timesteps].float().to(timesteps.device)
while len(res.shape) < len(broadcast_shape):
res = res[..., None]
return res.expand(broadcast_shape)
def parse_args():
parser = argparse.ArgumentParser(description="Simple example of a training script.")
parser.add_argument(
"--dataset_name",
type=str,
default=None,
help=(
"The name of the Dataset (from the HuggingFace hub) to train on (could be your own, possibly private,"
" dataset). It can also be a path pointing to a local copy of a dataset in your filesystem,"
" or to a folder containing files that HF Datasets can understand."
),
)
parser.add_argument(
"--dataset_config_name",
type=str,
default=None,
help="The config of the Dataset, leave as None if there's only one config.",
)
parser.add_argument(
"--train_data_dir",
type=str,
default=None,
help=(
"A folder containing the training data. Folder contents must follow the structure described in"
" https://huggingface.co/docs/datasets/image_dataset#imagefolder. In particular, a `metadata.jsonl` file"
" must exist to provide the captions for the images. Ignored if `dataset_name` is specified."
),
)
parser.add_argument(
"--output_dir",
type=str,
default="ddpm-model-64",
help="The output directory where the model predictions and checkpoints will be written.",
)
parser.add_argument("--overwrite_output_dir", action="store_true")
parser.add_argument(
"--cache_dir",
type=str,
default=None,
help="The directory where the downloaded models and datasets will be stored.",
)
parser.add_argument(
"--resolution",
type=int,
default=64,
help=(
"The resolution for input images, all the images in the train/validation dataset will be resized to this"
" resolution"
),
)
parser.add_argument(
"--train_batch_size", type=int, default=16, help="Batch size (per device) for the training dataloader."
)
parser.add_argument(
"--eval_batch_size", type=int, default=16, help="The number of images to generate for evaluation."
)
parser.add_argument(
"--dataloader_num_workers",
type=int,
default=0,
help=(
"The number of subprocesses to use for data loading. 0 means that the data will be loaded in the main"
" process."
),
)
parser.add_argument("--num_epochs", type=int, default=100)
parser.add_argument("--save_images_epochs", type=int, default=10, help="How often to save images during training.")
parser.add_argument(
"--save_model_epochs", type=int, default=10, help="How often to save the model during training."
)
parser.add_argument(
"--gradient_accumulation_steps",
type=int,
default=1,
help="Number of updates steps to accumulate before performing a backward/update pass.",
)
parser.add_argument(
"--learning_rate",
type=float,
default=1e-4,
help="Initial learning rate (after the potential warmup period) to use.",
)
parser.add_argument(
"--lr_scheduler",
type=str,
default="cosine",
help=(
'The scheduler type to use. Choose between ["linear", "cosine", "cosine_with_restarts", "polynomial",'
' "constant", "constant_with_warmup"]'
),
)
parser.add_argument(
"--lr_warmup_steps", type=int, default=500, help="Number of steps for the warmup in the lr scheduler."
)
parser.add_argument("--adam_beta1", type=float, default=0.95, help="The beta1 parameter for the Adam optimizer.")
parser.add_argument("--adam_beta2", type=float, default=0.999, help="The beta2 parameter for the Adam optimizer.")
parser.add_argument(
"--adam_weight_decay", type=float, default=1e-6, help="Weight decay magnitude for the Adam optimizer."
)
parser.add_argument("--adam_epsilon", type=float, default=1e-08, help="Epsilon value for the Adam optimizer.")
parser.add_argument(
"--use_ema",
action="store_true",
default=True,
help="Whether to use Exponential Moving Average for the final model weights.",
)
parser.add_argument("--ema_inv_gamma", type=float, default=1.0, help="The inverse gamma value for the EMA decay.")
parser.add_argument("--ema_power", type=float, default=3 / 4, help="The power value for the EMA decay.")
parser.add_argument("--ema_max_decay", type=float, default=0.9999, help="The maximum decay magnitude for EMA.")
parser.add_argument("--push_to_hub", action="store_true", help="Whether or not to push the model to the Hub.")
parser.add_argument("--hub_token", type=str, default=None, help="The token to use to push to the Model Hub.")
parser.add_argument(
"--hub_model_id",
type=str,
default=None,
help="The name of the repository to keep in sync with the local `output_dir`.",
)
parser.add_argument(
"--hub_private_repo", action="store_true", help="Whether or not to create a private repository."
)
parser.add_argument(
"--logging_dir",
type=str,
default="logs",
help=(
"[TensorBoard](https://www.tensorflow.org/tensorboard) log directory. Will default to"
" *output_dir/runs/**CURRENT_DATETIME_HOSTNAME***."
),
)
parser.add_argument("--local_rank", type=int, default=-1, help="For distributed training: local_rank")
parser.add_argument(
"--mixed_precision",
type=str,
default="no",
choices=["no", "fp16", "bf16"],
help=(
"Whether to use mixed precision. Choose"
"between fp16 and bf16 (bfloat16). Bf16 requires PyTorch >= 1.10."
"and an Nvidia Ampere GPU."
),
)
parser.add_argument(
"--prediction_type",
type=str,
default="epsilon",
choices=["epsilon", "sample"],
help="Whether the model should predict the 'epsilon'/noise error or directly the reconstructed image 'x0'.",
)
parser.add_argument("--ddpm_num_steps", type=int, default=1000)
parser.add_argument("--ddpm_beta_schedule", type=str, default="linear")
args = parser.parse_args()
env_local_rank = int(os.environ.get("LOCAL_RANK", -1))
if env_local_rank != -1 and env_local_rank != args.local_rank:
args.local_rank = env_local_rank
if args.dataset_name is None and args.train_data_dir is None:
raise ValueError("You must specify either a dataset name from the hub or a train data directory.")
return args
def get_full_repo_name(model_id: str, organization: Optional[str] = None, token: Optional[str] = None):
if token is None:
token = HfFolder.get_token()
if organization is None:
username = whoami(token)["name"]
return f"{username}/{model_id}"
else:
return f"{organization}/{model_id}"
def main(args):
logging_dir = os.path.join(args.output_dir, args.logging_dir)
accelerator = Accelerator(
@@ -61,7 +261,17 @@ def main(args):
),
)
model = ORTModule(model)
noise_scheduler = DDPMScheduler(num_train_timesteps=1000, tensor_format="pt")
accepts_prediction_type = "prediction_type" in set(inspect.signature(DDPMScheduler.__init__).parameters.keys())
if accepts_prediction_type:
noise_scheduler = DDPMScheduler(
num_train_timesteps=args.ddpm_num_steps,
beta_schedule=args.ddpm_beta_schedule,
prediction_type=args.prediction_type,
)
else:
noise_scheduler = DDPMScheduler(num_train_timesteps=args.ddpm_num_steps, beta_schedule=args.ddpm_beta_schedule)
optimizer = torch.optim.AdamW(
model.parameters(),
lr=args.learning_rate,
@@ -85,7 +295,6 @@ def main(args):
args.dataset_name,
args.dataset_config_name,
cache_dir=args.cache_dir,
use_auth_token=True if args.use_auth_token else None,
split="train",
)
else:
@@ -95,8 +304,12 @@ def main(args):
images = [augmentations(image.convert("RGB")) for image in examples["image"]]
return {"input": images}
logger.info(f"Dataset size: {len(dataset)}")
dataset.set_transform(transforms)
train_dataloader = torch.utils.data.DataLoader(dataset, batch_size=args.train_batch_size, shuffle=True)
train_dataloader = torch.utils.data.DataLoader(
dataset, batch_size=args.train_batch_size, shuffle=True, num_workers=args.dataloader_num_workers
)
lr_scheduler = get_scheduler(
args.lr_scheduler,
@@ -111,10 +324,29 @@ def main(args):
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
ema_model = EMAModel(model, inv_gamma=args.ema_inv_gamma, power=args.ema_power, max_value=args.ema_max_decay)
ema_model = EMAModel(
accelerator.unwrap_model(model),
inv_gamma=args.ema_inv_gamma,
power=args.ema_power,
max_value=args.ema_max_decay,
)
if args.push_to_hub:
repo = init_git_repo(args, at_init=True)
# Handle the repository creation
if accelerator.is_main_process:
if args.push_to_hub:
if args.hub_model_id is None:
repo_name = get_full_repo_name(Path(args.output_dir).name, token=args.hub_token)
else:
repo_name = args.hub_model_id
repo = Repository(args.output_dir, clone_from=repo_name)
with open(os.path.join(args.output_dir, ".gitignore"), "w+") as gitignore:
if "step_*" not in gitignore:
gitignore.write("step_*\n")
if "epoch_*" not in gitignore:
gitignore.write("epoch_*\n")
elif args.output_dir is not None:
os.makedirs(args.output_dir, exist_ok=True)
if accelerator.is_main_process:
run = os.path.split(__file__)[-1].split(".")[0]
@@ -141,11 +373,26 @@ def main(args):
with accelerator.accumulate(model):
# Predict the noise residual
noise_pred = model(noisy_images, timesteps, return_dict=True)[0]
loss = F.mse_loss(noise_pred, noise)
model_output = model(noisy_images, timesteps, return_dict=True)[0]
if args.prediction_type == "epsilon":
loss = F.mse_loss(model_output, noise) # this could have different weights!
elif args.prediction_type == "sample":
alpha_t = _extract_into_tensor(
noise_scheduler.alphas_cumprod, timesteps, (clean_images.shape[0], 1, 1, 1)
)
snr_weights = alpha_t / (1 - alpha_t)
loss = snr_weights * F.mse_loss(
model_output, clean_images, reduction="none"
) # use SNR weighting from distillation paper
loss = loss.mean()
else:
raise ValueError(f"Unsupported prediction type: {args.prediction_type}")
accelerator.backward(loss)
accelerator.clip_grad_norm_(model.parameters(), 1.0)
if accelerator.sync_gradients:
accelerator.clip_grad_norm_(model.parameters(), 1.0)
optimizer.step()
lr_scheduler.step()
if args.use_ema:
@@ -174,9 +421,13 @@ def main(args):
scheduler=noise_scheduler,
)
generator = torch.manual_seed(0)
generator = torch.Generator(device=pipeline.device).manual_seed(0)
# run pipeline in inference (sample random noise and denoise)
images = pipeline(generator=generator, batch_size=args.eval_batch_size, output_type="numpy").images
images = pipeline(
generator=generator,
batch_size=args.eval_batch_size,
output_type="numpy",
).images
# denormalize the images and save to tensorboard
images_processed = (images * 255).round().astype("uint8")
@@ -186,66 +437,14 @@ def main(args):
if epoch % args.save_model_epochs == 0 or epoch == args.num_epochs - 1:
# save the model
pipeline.save_pretrained(args.output_dir)
if args.push_to_hub:
push_to_hub(args, pipeline, repo, commit_message=f"Epoch {epoch}", blocking=False)
else:
pipeline.save_pretrained(args.output_dir)
repo.push_to_hub(commit_message=f"Epoch {epoch}", blocking=False)
accelerator.wait_for_everyone()
accelerator.end_training()
if __name__ == "__main__":
parser = argparse.ArgumentParser(description="Simple example of a training script.")
parser.add_argument("--local_rank", type=int, default=-1)
parser.add_argument("--dataset_name", type=str, default=None)
parser.add_argument("--dataset_config_name", type=str, default=None)
parser.add_argument("--train_data_dir", type=str, default=None, help="A folder containing the training data.")
parser.add_argument("--output_dir", type=str, default="ddpm-model-64")
parser.add_argument("--overwrite_output_dir", action="store_true")
parser.add_argument("--cache_dir", type=str, default=None)
parser.add_argument("--resolution", type=int, default=64)
parser.add_argument("--train_batch_size", type=int, default=16)
parser.add_argument("--eval_batch_size", type=int, default=16)
parser.add_argument("--num_epochs", type=int, default=100)
parser.add_argument("--save_images_epochs", type=int, default=10)
parser.add_argument("--save_model_epochs", type=int, default=10)
parser.add_argument("--gradient_accumulation_steps", type=int, default=1)
parser.add_argument("--learning_rate", type=float, default=1e-4)
parser.add_argument("--lr_scheduler", type=str, default="cosine")
parser.add_argument("--lr_warmup_steps", type=int, default=500)
parser.add_argument("--adam_beta1", type=float, default=0.95)
parser.add_argument("--adam_beta2", type=float, default=0.999)
parser.add_argument("--adam_weight_decay", type=float, default=1e-6)
parser.add_argument("--adam_epsilon", type=float, default=1e-08)
parser.add_argument("--use_ema", action="store_true", default=True)
parser.add_argument("--ema_inv_gamma", type=float, default=1.0)
parser.add_argument("--ema_power", type=float, default=3 / 4)
parser.add_argument("--ema_max_decay", type=float, default=0.9999)
parser.add_argument("--push_to_hub", action="store_true")
parser.add_argument("--use_auth_token", action="store_true")
parser.add_argument("--hub_token", type=str, default=None)
parser.add_argument("--hub_model_id", type=str, default=None)
parser.add_argument("--hub_private_repo", action="store_true")
parser.add_argument("--logging_dir", type=str, default="logs")
parser.add_argument(
"--mixed_precision",
type=str,
default="no",
choices=["no", "fp16", "bf16"],
help=(
"Whether to use mixed precision. Choose"
"between fp16 and bf16 (bfloat16). Bf16 requires PyTorch >= 1.10."
"and an Nvidia Ampere GPU."
),
)
args = parser.parse_args()
env_local_rank = int(os.environ.get("LOCAL_RANK", -1))
if env_local_rank != -1 and env_local_rank != args.local_rank:
args.local_rank = env_local_rank
if args.dataset_name is None and args.train_data_dir is None:
raise ValueError("You must specify either a dataset name from the hub or a train data directory.")
args = parse_args()
main(args)

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