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12 Commits

Author SHA1 Message Date
Dhruv Nair
be55fa631f update 2024-08-13 14:11:47 +02:00
Linoy Tsaban
413ca29b71 [Flux Dreambooth LoRA] - te bug fixes & updates (#9139)
* add requirements + fix link to bghira's guide

* text ecnoder training fixes

* text encoder training fixes

* text encoder training fixes

* text encoder training fixes

* style

* add tests

* fix encode_prompt call

* style

* unpack_latents test

* fix lora saving

* remove default val for max_sequenece_length in encode_prompt

* remove default val for max_sequenece_length in encode_prompt

* style

* testing

* style

* testing

* testing

* style

* fix sizing issue

* style

* revert scaling

* style

* style

* scaling test

* style

* scaling test

* remove model pred operation left from pre-conditioning

* remove model pred operation left from pre-conditioning

* fix trainable params

* remove te2 from casting

* transformer to accelerator

* remove prints

* empty commit
2024-08-12 11:58:03 +05:30
Dhruv Nair
10dc06c8d9 Update Video Loading/Export to use imageio (#9094)
* update

* update

* update

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-08-12 10:19:53 +05:30
Dibbla!
3ece143308 Errata - fix typo (#9100) 2024-08-12 07:30:19 +05:30
Steven Liu
98930ee131 [docs] Resolve internal links to PEFT (#9144)
* resolve peft links

* fuse_lora
2024-08-10 06:37:46 +05:30
Daniel Socek
c1079f0887 Fix textual inversion SDXL and add support for 2nd text encoder (#9010)
* Fix textual inversion SDXL and add support for 2nd text encoder

Signed-off-by: Daniel Socek <daniel.socek@intel.com>

* Fix style/quality of text inv for sdxl

Signed-off-by: Daniel Socek <daniel.socek@intel.com>

---------

Signed-off-by: Daniel Socek <daniel.socek@intel.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-08-09 20:23:06 +05:30
Linoy Tsaban
65e30907b5 [Flux] Dreambooth LoRA training scripts (#9086)
* initial commit - dreambooth for flux

* update transformer to be FluxTransformer2DModel

* update training loop and validation inference

* fix sd3->flux docs

* add guidance handling, not sure if it makes sense(?)

* inital dreambooth lora commit

* fix text_ids in compute_text_embeddings

* fix imports of static methods

* fix pipeline loading in readme, remove auto1111 docs for now

* fix pipeline loading in readme, remove auto1111 docs for now, remove some irrelevant text_encoder_3 refs

* Update examples/dreambooth/train_dreambooth_flux.py

Co-authored-by: Bagheera <59658056+bghira@users.noreply.github.com>

* fix te2 loading and remove te2 refs from text encoder training

* fix tokenizer_2 initialization

* remove text_encoder training refs from lora script (for now)

* try with vae in bfloat16, fix model hook save

* fix tokenization

* fix static imports

* fix CLIP import

* remove text_encoder training refs (for now) from lora script

* fix minor bug in encode_prompt, add guidance def in lora script, ...

* fix unpack_latents args

* fix license in readme

* add "none" to weighting_scheme options for uniform sampling

* style

* adapt model saving - remove text encoder refs

* adapt model loading - remove text encoder refs

* initial commit for readme

* Update examples/dreambooth/train_dreambooth_lora_flux.py

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* Update examples/dreambooth/train_dreambooth_lora_flux.py

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>

* fix vae casting

* remove precondition_outputs

* readme

* readme

* style

* readme

* readme

* update weighting scheme default & docs

* style

* add text_encoder training to lora script, change vae_scale_factor value in both

* style

* text encoder training fixes

* style

* update readme

* minor fixes

* fix te params

* fix te params

---------

Co-authored-by: Bagheera <59658056+bghira@users.noreply.github.com>
Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-08-09 07:31:04 +05:30
Sayak Paul
cee7c1b0fb Update README.md to include InstantID (#8770)
Co-authored-by: YiYi Xu <yixu310@gmail.com>
2024-08-08 10:14:12 -07:00
Monjoy Narayan Choudhury
1fcb811a8e Add Differential Diffusion to HunyuanDiT. (#9040)
* Add Differential Pipeline.

* Fix Styling Issue using ruff -fix

* Add details to Contributing.md

* Revert "Fix Styling Issue using ruff -fix"

This reverts commit d347de162d.

* Revert "Revert "Fix Styling Issue using ruff -fix""

This reverts commit ce7c3ff216.

* Revert README changes

* Restore README.md

* Update README.md

* Resolved Comments:

* Fix Readme based on review

* Fix formatting after make style

---------

Co-authored-by: Aryan <aryan@huggingface.co>
2024-08-08 18:53:39 +05:30
David Steinberg
ae026db7aa Fix a dead link (#9116)
Co-authored-by: Aryan <aryan@huggingface.co>
2024-08-08 18:46:50 +05:30
sayantan sadhu
8e3affc669 fix for lr scheduler in distributed training (#9103)
* fix for lr scheduler in distributed training

* Fixed the recalculation of the total training step section

* Fixed lint error

---------

Co-authored-by: Sayak Paul <spsayakpaul@gmail.com>
2024-08-08 08:45:48 +05:30
Steven Liu
ba7e48455a [docs] Organize model toctree (#9118)
* toctree

* fix
2024-08-08 08:31:58 +05:30
23 changed files with 5791 additions and 137 deletions

View File

@@ -202,6 +202,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
- https://github.com/microsoft/TaskMatrix
- https://github.com/invoke-ai/InvokeAI
- https://github.com/InstantID/InstantID
- https://github.com/apple/ml-stable-diffusion
- https://github.com/Sanster/lama-cleaner
- https://github.com/IDEA-Research/Grounded-Segment-Anything

View File

@@ -223,68 +223,76 @@
sections:
- local: api/models/overview
title: Overview
- local: api/models/unet
title: UNet1DModel
- local: api/models/unet2d
title: UNet2DModel
- local: api/models/unet2d-cond
title: UNet2DConditionModel
- local: api/models/unet3d-cond
title: UNet3DConditionModel
- local: api/models/unet-motion
title: UNetMotionModel
- local: api/models/uvit2d
title: UViT2DModel
- local: api/models/vq
title: VQModel
- local: api/models/autoencoderkl
title: AutoencoderKL
- local: api/models/autoencoderkl_cogvideox
title: AutoencoderKLCogVideoX
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/stable_cascade_unet
title: StableCascadeUNet
- local: api/models/autoencoder_tiny
title: Tiny AutoEncoder
- local: api/models/autoencoder_oobleck
title: Oobleck AutoEncoder
- local: api/models/consistency_decoder_vae
title: ConsistencyDecoderVAE
- local: api/models/transformer2d
title: Transformer2DModel
- local: api/models/pixart_transformer2d
title: PixArtTransformer2DModel
- local: api/models/dit_transformer2d
title: DiTTransformer2DModel
- local: api/models/hunyuan_transformer2d
title: HunyuanDiT2DModel
- local: api/models/aura_flow_transformer2d
title: AuraFlowTransformer2DModel
- local: api/models/flux_transformer
title: FluxTransformer2DModel
- local: api/models/latte_transformer3d
title: LatteTransformer3DModel
- local: api/models/cogvideox_transformer3d
title: CogVideoXTransformer3DModel
- local: api/models/lumina_nextdit2d
title: LuminaNextDiT2DModel
- local: api/models/transformer_temporal
title: TransformerTemporalModel
- local: api/models/sd3_transformer2d
title: SD3Transformer2DModel
- local: api/models/stable_audio_transformer
title: StableAudioDiTModel
- local: api/models/prior_transformer
title: PriorTransformer
- local: api/models/controlnet
title: ControlNetModel
- local: api/models/controlnet_hunyuandit
title: HunyuanDiT2DControlNetModel
- local: api/models/controlnet_sd3
title: SD3ControlNetModel
- local: api/models/controlnet_sparsectrl
title: SparseControlNetModel
- sections:
- local: api/models/controlnet
title: ControlNetModel
- local: api/models/controlnet_hunyuandit
title: HunyuanDiT2DControlNetModel
- local: api/models/controlnet_sd3
title: SD3ControlNetModel
- local: api/models/controlnet_sparsectrl
title: SparseControlNetModel
title: ControlNets
- sections:
- local: api/models/aura_flow_transformer2d
title: AuraFlowTransformer2DModel
- local: api/models/cogvideox_transformer3d
title: CogVideoXTransformer3DModel
- local: api/models/dit_transformer2d
title: DiTTransformer2DModel
- local: api/models/flux_transformer
title: FluxTransformer2DModel
- local: api/models/hunyuan_transformer2d
title: HunyuanDiT2DModel
- local: api/models/latte_transformer3d
title: LatteTransformer3DModel
- local: api/models/lumina_nextdit2d
title: LuminaNextDiT2DModel
- local: api/models/pixart_transformer2d
title: PixArtTransformer2DModel
- local: api/models/prior_transformer
title: PriorTransformer
- local: api/models/sd3_transformer2d
title: SD3Transformer2DModel
- local: api/models/stable_audio_transformer
title: StableAudioDiTModel
- local: api/models/transformer2d
title: Transformer2DModel
- local: api/models/transformer_temporal
title: TransformerTemporalModel
title: Transformers
- sections:
- local: api/models/stable_cascade_unet
title: StableCascadeUNet
- local: api/models/unet
title: UNet1DModel
- local: api/models/unet2d
title: UNet2DModel
- local: api/models/unet2d-cond
title: UNet2DConditionModel
- local: api/models/unet3d-cond
title: UNet3DConditionModel
- local: api/models/unet-motion
title: UNetMotionModel
- local: api/models/uvit2d
title: UViT2DModel
title: UNets
- sections:
- local: api/models/autoencoderkl
title: AutoencoderKL
- local: api/models/autoencoderkl_cogvideox
title: AutoencoderKLCogVideoX
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/consistency_decoder_vae
title: ConsistencyDecoderVAE
- local: api/models/autoencoder_oobleck
title: Oobleck AutoEncoder
- local: api/models/autoencoder_tiny
title: Tiny AutoEncoder
- local: api/models/vq
title: VQModel
title: VAEs
title: Models
- isExpanded: false
sections:

View File

@@ -21,7 +21,7 @@ Stable Audio is trained on a corpus of around 48k audio recordings, where around
The abstract of the paper is the following:
*Open generative models are vitally important for the community, allowing for fine-tunes and serving as baselines when presenting new models. However, most current text-to-audio models are private and not accessible for artists and researchers to build upon. Here we describe the architecture and training process of a new open-weights text-to-audio model trained with Creative Commons data. Our evaluation shows that the model's performance is competitive with the state-of-the-art across various metrics. Notably, the reported FDopenl3 results (measuring the realism of the generations) showcase its potential for high-quality stereo sound synthesis at 44.1kHz.*
This pipeline was contributed by [Yoach Lacombe](https://huggingface.co/ylacombe). The original codebase can be found at [Stability-AI/stable-audio-tool](https://github.com/Stability-AI/stable-audio-tool).
This pipeline was contributed by [Yoach Lacombe](https://huggingface.co/ylacombe). The original codebase can be found at [Stability-AI/stable-audio-tools](https://github.com/Stability-AI/stable-audio-tools).
## Tips

View File

@@ -125,3 +125,5 @@ image
<figcaption class="mt-2 text-center text-sm text-gray-500">distilled Stable Diffusion + Tiny AutoEncoder</figcaption>
</div>
</div>
More tiny autoencoder models for other Stable Diffusion models, like Stable Diffusion 3, are available from [madebyollin](https://huggingface.co/madebyollin).

View File

@@ -14,9 +14,9 @@ specific language governing permissions and limitations under the License.
It can be fun and creative to use multiple [LoRAs]((https://huggingface.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora)) together to generate something entirely new and unique. This works by merging multiple LoRA weights together to produce images that are a blend of different styles. Diffusers provides a few methods to merge LoRAs depending on *how* you want to merge their weights, which can affect image quality.
This guide will show you how to merge LoRAs using the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] and [`~peft.LoraModel.add_weighted_adapter`] methods. To improve inference speed and reduce memory-usage of merged LoRAs, you'll also see how to use the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method to fuse the LoRA weights with the original weights of the underlying model.
This guide will show you how to merge LoRAs using the [`~loaders.PeftAdapterMixin.set_adapters`] and [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) methods. To improve inference speed and reduce memory-usage of merged LoRAs, you'll also see how to use the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method to fuse the LoRA weights with the original weights of the underlying model.
For this guide, load a Stable Diffusion XL (SDXL) checkpoint and the [KappaNeuro/studio-ghibli-style]() and [Norod78/sdxl-chalkboarddrawing-lora]() LoRAs with the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method. You'll need to assign each LoRA an `adapter_name` to combine them later.
For this guide, load a Stable Diffusion XL (SDXL) checkpoint and the [KappaNeuro/studio-ghibli-style](https://huggingface.co/KappaNeuro/studio-ghibli-style) and [Norod78/sdxl-chalkboarddrawing-lora](https://huggingface.co/Norod78/sdxl-chalkboarddrawing-lora) LoRAs with the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] method. You'll need to assign each LoRA an `adapter_name` to combine them later.
```py
from diffusers import DiffusionPipeline
@@ -29,7 +29,7 @@ pipeline.load_lora_weights("lordjia/by-feng-zikai", weight_name="fengzikai_v1.0_
## set_adapters
The [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] method merges LoRA adapters by concatenating their weighted matrices. Use the adapter name to specify which LoRAs to merge, and the `adapter_weights` parameter to control the scaling for each LoRA. For example, if `adapter_weights=[0.5, 0.5]`, then the merged LoRA output is an average of both LoRAs. Try adjusting the adapter weights to see how it affects the generated image!
The [`~loaders.PeftAdapterMixin.set_adapters`] method merges LoRA adapters by concatenating their weighted matrices. Use the adapter name to specify which LoRAs to merge, and the `adapter_weights` parameter to control the scaling for each LoRA. For example, if `adapter_weights=[0.5, 0.5]`, then the merged LoRA output is an average of both LoRAs. Try adjusting the adapter weights to see how it affects the generated image!
```py
pipeline.set_adapters(["ikea", "feng"], adapter_weights=[0.7, 0.8])
@@ -47,19 +47,19 @@ image
## add_weighted_adapter
> [!WARNING]
> This is an experimental method that adds PEFTs [`~peft.LoraModel.add_weighted_adapter`] method to Diffusers to enable more efficient merging methods. Check out this [issue](https://github.com/huggingface/diffusers/issues/6892) if you're interested in learning more about the motivation and design behind this integration.
> This is an experimental method that adds PEFTs [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) method to Diffusers to enable more efficient merging methods. Check out this [issue](https://github.com/huggingface/diffusers/issues/6892) if you're interested in learning more about the motivation and design behind this integration.
The [`~peft.LoraModel.add_weighted_adapter`] method provides access to more efficient merging method such as [TIES and DARE](https://huggingface.co/docs/peft/developer_guides/model_merging). To use these merging methods, make sure you have the latest stable version of Diffusers and PEFT installed.
The [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) method provides access to more efficient merging method such as [TIES and DARE](https://huggingface.co/docs/peft/developer_guides/model_merging). To use these merging methods, make sure you have the latest stable version of Diffusers and PEFT installed.
```bash
pip install -U diffusers peft
```
There are three steps to merge LoRAs with the [`~peft.LoraModel.add_weighted_adapter`] method:
There are three steps to merge LoRAs with the [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) method:
1. Create a [`~peft.PeftModel`] from the underlying model and LoRA checkpoint.
1. Create a [PeftModel](https://huggingface.co/docs/peft/package_reference/peft_model#peft.PeftModel) from the underlying model and LoRA checkpoint.
2. Load a base UNet model and the LoRA adapters.
3. Merge the adapters using the [`~peft.LoraModel.add_weighted_adapter`] method and the merging method of your choice.
3. Merge the adapters using the [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) method and the merging method of your choice.
Let's dive deeper into what these steps entail.
@@ -92,7 +92,7 @@ pipeline = DiffusionPipeline.from_pretrained(
pipeline.load_lora_weights("ostris/ikea-instructions-lora-sdxl", weight_name="ikea_instructions_xl_v1_5.safetensors", adapter_name="ikea")
```
Now you'll create a [`~peft.PeftModel`] from the loaded LoRA checkpoint by combining the SDXL UNet and the LoRA UNet from the pipeline.
Now you'll create a [PeftModel](https://huggingface.co/docs/peft/package_reference/peft_model#peft.PeftModel) from the loaded LoRA checkpoint by combining the SDXL UNet and the LoRA UNet from the pipeline.
```python
from peft import get_peft_model, LoraConfig
@@ -112,7 +112,7 @@ ikea_peft_model.load_state_dict(original_state_dict, strict=True)
> [!TIP]
> You can optionally push the ikea_peft_model to the Hub by calling `ikea_peft_model.push_to_hub("ikea_peft_model", token=TOKEN)`.
Repeat this process to create a [`~peft.PeftModel`] from the [lordjia/by-feng-zikai](https://huggingface.co/lordjia/by-feng-zikai) LoRA.
Repeat this process to create a [PeftModel](https://huggingface.co/docs/peft/package_reference/peft_model#peft.PeftModel) from the [lordjia/by-feng-zikai](https://huggingface.co/lordjia/by-feng-zikai) LoRA.
```python
pipeline.delete_adapters("ikea")
@@ -148,7 +148,7 @@ model = PeftModel.from_pretrained(base_unet, "stevhliu/ikea_peft_model", use_saf
model.load_adapter("stevhliu/feng_peft_model", use_safetensors=True, subfolder="feng", adapter_name="feng")
```
3. Merge the adapters using the [`~peft.LoraModel.add_weighted_adapter`] method and the merging method of your choice (learn more about other merging methods in this [blog post](https://huggingface.co/blog/peft_merging)). For this example, let's use the `"dare_linear"` method to merge the LoRAs.
3. Merge the adapters using the [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) method and the merging method of your choice (learn more about other merging methods in this [blog post](https://huggingface.co/blog/peft_merging)). For this example, let's use the `"dare_linear"` method to merge the LoRAs.
> [!WARNING]
> Keep in mind the LoRAs need to have the same rank to be merged!
@@ -182,9 +182,9 @@ image
## fuse_lora
Both the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] and [`~peft.LoraModel.add_weighted_adapter`] methods require loading the base model and the LoRA adapters separately which incurs some overhead. The [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method allows you to fuse the LoRA weights directly with the original weights of the underlying model. This way, you're only loading the model once which can increase inference and lower memory-usage.
Both the [`~loaders.PeftAdapterMixin.set_adapters`] and [add_weighted_adapter](https://huggingface.co/docs/peft/package_reference/lora#peft.LoraModel.add_weighted_adapter) methods require loading the base model and the LoRA adapters separately which incurs some overhead. The [`~loaders.lora_base.LoraBaseMixin.fuse_lora`] method allows you to fuse the LoRA weights directly with the original weights of the underlying model. This way, you're only loading the model once which can increase inference and lower memory-usage.
You can use PEFT to easily fuse/unfuse multiple adapters directly into the model weights (both UNet and text encoder) using the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method, which can lead to a speed-up in inference and lower VRAM usage.
You can use PEFT to easily fuse/unfuse multiple adapters directly into the model weights (both UNet and text encoder) using the [`~loaders.lora_base.LoraBaseMixin.fuse_lora`] method, which can lead to a speed-up in inference and lower VRAM usage.
For example, if you have a base model and adapters loaded and set as active with the following adapter weights:
@@ -199,7 +199,7 @@ pipeline.load_lora_weights("lordjia/by-feng-zikai", weight_name="fengzikai_v1.0_
pipeline.set_adapters(["ikea", "feng"], adapter_weights=[0.7, 0.8])
```
Fuse these LoRAs into the UNet with the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method. The `lora_scale` parameter controls how much to scale the output by with the LoRA weights. It is important to make the `lora_scale` adjustments in the [`~loaders.StableDiffusionLoraLoaderMixin.fuse_lora`] method because it wont work if you try to pass `scale` to the `cross_attention_kwargs` in the pipeline.
Fuse these LoRAs into the UNet with the [`~loaders.lora_base.LoraBaseMixin.fuse_lora`] method. The `lora_scale` parameter controls how much to scale the output by with the LoRA weights. It is important to make the `lora_scale` adjustments in the [`~loaders.lora_base.LoraBaseMixin.fuse_lora`] method because it wont work if you try to pass `scale` to the `cross_attention_kwargs` in the pipeline.
```py
pipeline.fuse_lora(adapter_names=["ikea", "feng"], lora_scale=1.0)
@@ -226,7 +226,7 @@ image = pipeline("A bowl of ramen shaped like a cute kawaii bear, by Feng Zikai"
image
```
You can call [`~loaders.StableDiffusionLoraLoaderMixin.unfuse_lora`] to restore the original model's weights (for example, if you want to use a different `lora_scale` value). However, this only works if you've only fused one LoRA adapter to the original model. If you've fused multiple LoRAs, you'll need to reload the model.
You can call [`~~loaders.lora_base.LoraBaseMixin.unfuse_lora`] to restore the original model's weights (for example, if you want to use a different `lora_scale` value). However, this only works if you've only fused one LoRA adapter to the original model. If you've fused multiple LoRAs, you'll need to reload the model.
```py
pipeline.unfuse_lora()

View File

@@ -71,6 +71,7 @@ Please also check out our [Community Scripts](https://github.com/huggingface/dif
| Stable Diffusion BoxDiff Pipeline | Training-free controlled generation with bounding boxes using [BoxDiff](https://github.com/showlab/BoxDiff) | [Stable Diffusion BoxDiff Pipeline](#stable-diffusion-boxdiff) | - | [Jingyang Zhang](https://github.com/zjysteven/) |
| FRESCO V2V Pipeline | Implementation of [[CVPR 2024] FRESCO: Spatial-Temporal Correspondence for Zero-Shot Video Translation](https://arxiv.org/abs/2403.12962) | [FRESCO V2V Pipeline](#fresco) | - | [Yifan Zhou](https://github.com/SingleZombie) |
| AnimateDiff IPEX Pipeline | Accelerate AnimateDiff inference pipeline with BF16/FP32 precision on Intel Xeon CPUs with [IPEX](https://github.com/intel/intel-extension-for-pytorch) | [AnimateDiff on IPEX](#animatediff-on-ipex) | - | [Dan Li](https://github.com/ustcuna/) |
| HunyuanDiT Differential Diffusion Pipeline | Applies [Differential Diffsuion](https://github.com/exx8/differential-diffusion) to [HunyuanDiT](https://github.com/huggingface/diffusers/pull/8240). | [HunyuanDiT with Differential Diffusion](#hunyuandit-with-differential-diffusion) | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/drive/1v44a5fpzyr4Ffr4v2XBQ7BajzG874N4P?usp=sharing) | [Monjoy Choudhury](https://github.com/MnCSSJ4x) |
To load a custom pipeline you just need to pass the `custom_pipeline` argument to `DiffusionPipeline`, as one of the files in `diffusers/examples/community`. Feel free to send a PR with your own pipelines, we will merge them quickly.
@@ -1646,7 +1647,6 @@ from diffusers import DiffusionPipeline
scheduler = DDIMScheduler.from_pretrained("stabilityai/stable-diffusion-2-1",
subfolder="scheduler")
pipe = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1",
custom_pipeline="stable_diffusion_tensorrt_img2img",
variant='fp16',
@@ -1661,7 +1661,6 @@ pipe = pipe.to("cuda")
url = "https://pajoca.com/wp-content/uploads/2022/09/tekito-yamakawa-1.png"
response = requests.get(url)
input_image = Image.open(BytesIO(response.content)).convert("RGB")
prompt = "photorealistic new zealand hills"
image = pipe(prompt, image=input_image, strength=0.75,).images[0]
image.save('tensorrt_img2img_new_zealand_hills.png')
@@ -4209,6 +4208,52 @@ print("Latency of AnimateDiffPipelineIpex--fp32", latency, "s for total", step,
latency = elapsed_time(pipe4, num_inference_steps=step)
print("Latency of AnimateDiffPipeline--fp32",latency, "s for total", step, "steps")
```
### HunyuanDiT with Differential Diffusion
#### Usage
```python
import torch
from diffusers import FlowMatchEulerDiscreteScheduler
from diffusers.utils import load_image
from PIL import Image
from torchvision import transforms
from pipeline_hunyuandit_differential_img2img import (
HunyuanDiTDifferentialImg2ImgPipeline,
)
pipe = HunyuanDiTDifferentialImg2ImgPipeline.from_pretrained(
"Tencent-Hunyuan/HunyuanDiT-Diffusers", torch_dtype=torch.float16
).to("cuda")
source_image = load_image(
"https://huggingface.co/datasets/OzzyGT/testing-resources/resolve/main/differential/20240329211129_4024911930.png"
)
map = load_image(
"https://huggingface.co/datasets/OzzyGT/testing-resources/resolve/main/differential/gradient_mask_2.png"
)
prompt = "a green pear"
negative_prompt = "blurry"
image = pipe(
prompt=prompt,
negative_prompt=negative_prompt,
image=source_image,
num_inference_steps=28,
guidance_scale=4.5,
strength=1.0,
map=map,
).images[0]
```
| ![Gradient](https://github.com/user-attachments/assets/e38ce4d5-1ae6-4df0-ab43-adc1b45716b5) | ![Input](https://github.com/user-attachments/assets/9c95679c-e9d7-4f5a-90d6-560203acd6b3) | ![Output](https://github.com/user-attachments/assets/5313ff64-a0c4-418b-8b55-a38f1a5e7532) |
| ------------------------------------------------------------------------------------------ | --------------------------------------------------------------------------------------- | ---------------------------------------------------------------------------------------- |
| Gradient | Input | Output |
A colab notebook demonstrating all results can be found [here](https://colab.research.google.com/drive/1v44a5fpzyr4Ffr4v2XBQ7BajzG874N4P?usp=sharing). Depth Maps have also been added in the same colab.
# Perturbed-Attention Guidance
@@ -4285,4 +4330,4 @@ grid_image.save(grid_dir + "sample.png")
`pag_scale` : guidance scale of PAG (ex: 5.0)
`pag_applied_layers_index` : index of the layer to apply perturbation (ex: ['m0'])
`pag_applied_layers_index` : index of the layer to apply perturbation (ex: ['m0'])

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# DreamBooth training example for FLUX.1 [dev]
[DreamBooth](https://arxiv.org/abs/2208.12242) is a method to personalize text2image models like stable diffusion given just a few (3~5) images of a subject.
The `train_dreambooth_flux.py` script shows how to implement the training procedure and adapt it for [FLUX.1 [dev]](https://blackforestlabs.ai/announcing-black-forest-labs/). We also provide a LoRA implementation in the `train_dreambooth_lora_flux.py` script.
> [!NOTE]
> **Memory consumption**
>
> Flux can be quite expensive to run on consumer hardware devices and as a result finetuning it comes with high memory requirements -
> a LoRA with a rank of 16 (w/ all components trained) can exceed 40GB of VRAM for training.
> For more tips & guidance on training on a resource-constrained device please visit [`@bghira`'s guide](https://github.com/bghira/SimpleTuner/blob/main/documentation/quickstart/FLUX.md)
> [!NOTE]
> **Gated model**
>
> As the model is gated, before using it with diffusers you first need to go to the [FLUX.1 [dev] Hugging Face page](https://huggingface.co/black-forest-labs/FLUX.1-dev), fill in the form and accept the gate. Once you are in, you need to log in so that your system knows youve accepted the gate. Use the command below to log in:
```bash
huggingface-cli login
```
This will also allow us to push the trained model parameters to the Hugging Face Hub platform.
## Running locally with PyTorch
### Installing the dependencies
Before running the scripts, make sure to install the library's training dependencies:
**Important**
To make sure you can successfully run the latest versions of the example scripts, we highly recommend **installing from source** and keeping the install up to date as we update the example scripts frequently and install some example-specific requirements. To do this, execute the following steps in a new virtual environment:
```bash
git clone https://github.com/huggingface/diffusers
cd diffusers
pip install -e .
```
Then cd in the `examples/dreambooth` folder and run
```bash
pip install -r requirements_flux.txt
```
And initialize an [🤗Accelerate](https://github.com/huggingface/accelerate/) environment with:
```bash
accelerate config
```
Or for a default accelerate configuration without answering questions about your environment
```bash
accelerate config default
```
Or if your environment doesn't support an interactive shell (e.g., a notebook)
```python
from accelerate.utils import write_basic_config
write_basic_config()
```
When running `accelerate config`, if we specify torch compile mode to True there can be dramatic speedups.
Note also that we use PEFT library as backend for LoRA training, make sure to have `peft>=0.6.0` installed in your environment.
### Dog toy example
Now let's get our dataset. For this example we will use some dog images: https://huggingface.co/datasets/diffusers/dog-example.
Let's first download it locally:
```python
from huggingface_hub import snapshot_download
local_dir = "./dog"
snapshot_download(
"diffusers/dog-example",
local_dir=local_dir, repo_type="dataset",
ignore_patterns=".gitattributes",
)
```
This will also allow us to push the trained LoRA parameters to the Hugging Face Hub platform.
Now, we can launch training using:
```bash
export MODEL_NAME="black-forest-labs/FLUX.1-dev"
export INSTANCE_DIR="dog"
export OUTPUT_DIR="trained-flux"
accelerate launch train_dreambooth_flux.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--output_dir=$OUTPUT_DIR \
--mixed_precision="bf16" \
--instance_prompt="a photo of sks dog" \
--resolution=1024 \
--train_batch_size=1 \
--gradient_accumulation_steps=4 \
--learning_rate=1e-4 \
--report_to="wandb" \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--max_train_steps=500 \
--validation_prompt="A photo of sks dog in a bucket" \
--validation_epochs=25 \
--seed="0" \
--push_to_hub
```
To better track our training experiments, we're using the following flags in the command above:
* `report_to="wandb` will ensure the training runs are tracked on Weights and Biases. To use it, be sure to install `wandb` with `pip install wandb`.
* `validation_prompt` and `validation_epochs` to allow the script to do a few validation inference runs. This allows us to qualitatively check if the training is progressing as expected.
> [!NOTE]
> If you want to train using long prompts with the T5 text encoder, you can use `--max_sequence_length` to set the token limit. The default is 77, but it can be increased to as high as 512. Note that this will use more resources and may slow down the training in some cases.
> [!TIP]
> You can pass `--use_8bit_adam` to reduce the memory requirements of training. Make sure to install `bitsandbytes` if you want to do so.
## LoRA + DreamBooth
[LoRA](https://huggingface.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora) is a popular parameter-efficient fine-tuning technique that allows you to achieve full-finetuning like performance but with a fraction of learnable parameters.
Note also that we use PEFT library as backend for LoRA training, make sure to have `peft>=0.6.0` installed in your environment.
To perform DreamBooth with LoRA, run:
```bash
export MODEL_NAME="black-forest-labs/FLUX.1-dev"
export INSTANCE_DIR="dog"
export OUTPUT_DIR="trained-flux-lora"
accelerate launch train_dreambooth_lora_flux.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--output_dir=$OUTPUT_DIR \
--mixed_precision="bf16" \
--instance_prompt="a photo of sks dog" \
--resolution=512 \
--train_batch_size=1 \
--gradient_accumulation_steps=4 \
--learning_rate=1e-5 \
--report_to="wandb" \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--max_train_steps=500 \
--validation_prompt="A photo of sks dog in a bucket" \
--validation_epochs=25 \
--seed="0" \
--push_to_hub
```
### Text Encoder Training
Alongside the transformer, fine-tuning of the CLIP text encoder is also supported.
To do so, just specify `--train_text_encoder` while launching training. Please keep the following points in mind:
> [!NOTE]
> FLUX.1 has 2 text encoders (CLIP L/14 and T5-v1.1-XXL).
By enabling `--train_text_encoder`, fine-tuning of the **CLIP encoder** is performed.
> At the moment, T5 fine-tuning is not supported and weights remain frozen when text encoder training is enabled.
To perform DreamBooth LoRA with text-encoder training, run:
```bash
export MODEL_NAME="black-forest-labs/FLUX.1-dev"
export OUTPUT_DIR="trained-flux-dev-dreambooth-lora"
accelerate launch train_dreambooth_lora_flux.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--output_dir=$OUTPUT_DIR \
--mixed_precision="bf16" \
--train_text_encoder\
--instance_prompt="a photo of sks dog" \
--resolution=512 \
--train_batch_size=1 \
--gradient_accumulation_steps=4 \
--learning_rate=1e-5 \
--report_to="wandb" \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--max_train_steps=500 \
--validation_prompt="A photo of sks dog in a bucket" \
--seed="0" \
--push_to_hub
```
## Other notes
Thanks to `bghira` for their help with reviewing & insight sharing ♥️

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@@ -0,0 +1,8 @@
accelerate>=0.31.0
torchvision
transformers>=4.41.2
ftfy
tensorboard
Jinja2
peft>=0.11.1
sentencepiece

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@@ -0,0 +1,203 @@
# coding=utf-8
# Copyright 2024 HuggingFace Inc.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import logging
import os
import shutil
import sys
import tempfile
from diffusers import DiffusionPipeline, FluxTransformer2DModel
sys.path.append("..")
from test_examples_utils import ExamplesTestsAccelerate, run_command # noqa: E402
logging.basicConfig(level=logging.DEBUG)
logger = logging.getLogger()
stream_handler = logging.StreamHandler(sys.stdout)
logger.addHandler(stream_handler)
class DreamBoothFlux(ExamplesTestsAccelerate):
instance_data_dir = "docs/source/en/imgs"
instance_prompt = "photo"
pretrained_model_name_or_path = "hf-internal-testing/tiny-flux-pipe"
script_path = "examples/dreambooth/train_dreambooth_flux.py"
def test_dreambooth(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
{self.script_path}
--pretrained_model_name_or_path {self.pretrained_model_name_or_path}
--instance_data_dir {self.instance_data_dir}
--instance_prompt {self.instance_prompt}
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "transformer", "diffusion_pytorch_model.safetensors")))
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "scheduler", "scheduler_config.json")))
def test_dreambooth_checkpointing(self):
with tempfile.TemporaryDirectory() as tmpdir:
# Run training script with checkpointing
# max_train_steps == 4, checkpointing_steps == 2
# Should create checkpoints at steps 2, 4
initial_run_args = f"""
{self.script_path}
--pretrained_model_name_or_path {self.pretrained_model_name_or_path}
--instance_data_dir {self.instance_data_dir}
--instance_prompt {self.instance_prompt}
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 4
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
--checkpointing_steps=2
--seed=0
""".split()
run_command(self._launch_args + initial_run_args)
# check can run the original fully trained output pipeline
pipe = DiffusionPipeline.from_pretrained(tmpdir)
pipe(self.instance_prompt, num_inference_steps=1)
# check checkpoint directories exist
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-2")))
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-4")))
# check can run an intermediate checkpoint
transformer = FluxTransformer2DModel.from_pretrained(tmpdir, subfolder="checkpoint-2/transformer")
pipe = DiffusionPipeline.from_pretrained(self.pretrained_model_name_or_path, transformer=transformer)
pipe(self.instance_prompt, num_inference_steps=1)
# Remove checkpoint 2 so that we can check only later checkpoints exist after resuming
shutil.rmtree(os.path.join(tmpdir, "checkpoint-2"))
# Run training script for 7 total steps resuming from checkpoint 4
resume_run_args = f"""
{self.script_path}
--pretrained_model_name_or_path {self.pretrained_model_name_or_path}
--instance_data_dir {self.instance_data_dir}
--instance_prompt {self.instance_prompt}
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 6
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
--checkpointing_steps=2
--resume_from_checkpoint=checkpoint-4
--seed=0
""".split()
run_command(self._launch_args + resume_run_args)
# check can run new fully trained pipeline
pipe = DiffusionPipeline.from_pretrained(tmpdir)
pipe(self.instance_prompt, num_inference_steps=1)
# check old checkpoints do not exist
self.assertFalse(os.path.isdir(os.path.join(tmpdir, "checkpoint-2")))
# check new checkpoints exist
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-4")))
self.assertTrue(os.path.isdir(os.path.join(tmpdir, "checkpoint-6")))
def test_dreambooth_checkpointing_checkpoints_total_limit(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
{self.script_path}
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
--instance_data_dir={self.instance_data_dir}
--output_dir={tmpdir}
--instance_prompt={self.instance_prompt}
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=6
--checkpoints_total_limit=2
--checkpointing_steps=2
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-4", "checkpoint-6"},
)
def test_dreambooth_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
{self.script_path}
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
--instance_data_dir={self.instance_data_dir}
--output_dir={tmpdir}
--instance_prompt={self.instance_prompt}
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=4
--checkpointing_steps=2
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-2", "checkpoint-4"},
)
resume_run_args = f"""
{self.script_path}
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
--instance_data_dir={self.instance_data_dir}
--output_dir={tmpdir}
--instance_prompt={self.instance_prompt}
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=8
--checkpointing_steps=2
--resume_from_checkpoint=checkpoint-4
--checkpoints_total_limit=2
""".split()
run_command(self._launch_args + resume_run_args)
self.assertEqual({x for x in os.listdir(tmpdir) if "checkpoint" in x}, {"checkpoint-6", "checkpoint-8"})

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@@ -0,0 +1,165 @@
# coding=utf-8
# Copyright 2024 HuggingFace Inc.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
import logging
import os
import sys
import tempfile
import safetensors
sys.path.append("..")
from test_examples_utils import ExamplesTestsAccelerate, run_command # noqa: E402
logging.basicConfig(level=logging.DEBUG)
logger = logging.getLogger()
stream_handler = logging.StreamHandler(sys.stdout)
logger.addHandler(stream_handler)
class DreamBoothLoRAFlux(ExamplesTestsAccelerate):
instance_data_dir = "docs/source/en/imgs"
instance_prompt = "photo"
pretrained_model_name_or_path = "hf-internal-testing/tiny-flux-pipe"
script_path = "examples/dreambooth/train_dreambooth_lora_flux.py"
def test_dreambooth_lora_flux(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
{self.script_path}
--pretrained_model_name_or_path {self.pretrained_model_name_or_path}
--instance_data_dir {self.instance_data_dir}
--instance_prompt {self.instance_prompt}
--resolution 64
--train_batch_size 1
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "pytorch_lora_weights.safetensors")))
# make sure the state_dict has the correct naming in the parameters.
lora_state_dict = safetensors.torch.load_file(os.path.join(tmpdir, "pytorch_lora_weights.safetensors"))
is_lora = all("lora" in k for k in lora_state_dict.keys())
self.assertTrue(is_lora)
# when not training the text encoder, all the parameters in the state dict should start
# with `"transformer"` in their names.
starts_with_transformer = all(key.startswith("transformer") for key in lora_state_dict.keys())
self.assertTrue(starts_with_transformer)
def test_dreambooth_lora_text_encoder_flux(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
{self.script_path}
--pretrained_model_name_or_path {self.pretrained_model_name_or_path}
--instance_data_dir {self.instance_data_dir}
--instance_prompt {self.instance_prompt}
--resolution 64
--train_batch_size 1
--train_text_encoder
--gradient_accumulation_steps 1
--max_train_steps 2
--learning_rate 5.0e-04
--scale_lr
--lr_scheduler constant
--lr_warmup_steps 0
--output_dir {tmpdir}
""".split()
run_command(self._launch_args + test_args)
# save_pretrained smoke test
self.assertTrue(os.path.isfile(os.path.join(tmpdir, "pytorch_lora_weights.safetensors")))
# make sure the state_dict has the correct naming in the parameters.
lora_state_dict = safetensors.torch.load_file(os.path.join(tmpdir, "pytorch_lora_weights.safetensors"))
is_lora = all("lora" in k for k in lora_state_dict.keys())
self.assertTrue(is_lora)
starts_with_expected_prefix = all(
(key.startswith("transformer") or key.startswith("text_encoder")) for key in lora_state_dict.keys()
)
self.assertTrue(starts_with_expected_prefix)
def test_dreambooth_lora_flux_checkpointing_checkpoints_total_limit(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
{self.script_path}
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
--instance_data_dir={self.instance_data_dir}
--output_dir={tmpdir}
--instance_prompt={self.instance_prompt}
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=6
--checkpoints_total_limit=2
--checkpointing_steps=2
""".split()
run_command(self._launch_args + test_args)
self.assertEqual(
{x for x in os.listdir(tmpdir) if "checkpoint" in x},
{"checkpoint-4", "checkpoint-6"},
)
def test_dreambooth_lora_flux_checkpointing_checkpoints_total_limit_removes_multiple_checkpoints(self):
with tempfile.TemporaryDirectory() as tmpdir:
test_args = f"""
{self.script_path}
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
--instance_data_dir={self.instance_data_dir}
--output_dir={tmpdir}
--instance_prompt={self.instance_prompt}
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=4
--checkpointing_steps=2
""".split()
run_command(self._launch_args + test_args)
self.assertEqual({x for x in os.listdir(tmpdir) if "checkpoint" in x}, {"checkpoint-2", "checkpoint-4"})
resume_run_args = f"""
{self.script_path}
--pretrained_model_name_or_path={self.pretrained_model_name_or_path}
--instance_data_dir={self.instance_data_dir}
--output_dir={tmpdir}
--instance_prompt={self.instance_prompt}
--resolution=64
--train_batch_size=1
--gradient_accumulation_steps=1
--max_train_steps=8
--checkpointing_steps=2
--resume_from_checkpoint=checkpoint-4
--checkpoints_total_limit=2
""".split()
run_command(self._launch_args + resume_run_args)
self.assertEqual({x for x in os.listdir(tmpdir) if "checkpoint" in x}, {"checkpoint-6", "checkpoint-8"})

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@@ -826,17 +826,22 @@ def main():
)
# Scheduler and math around the number of training steps.
overrode_max_train_steps = False
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
# Check the PR https://github.com/huggingface/diffusers/pull/8312 for detailed explanation.
num_warmup_steps_for_scheduler = args.lr_warmup_steps * accelerator.num_processes
if args.max_train_steps is None:
args.max_train_steps = args.num_train_epochs * num_update_steps_per_epoch
overrode_max_train_steps = True
len_train_dataloader_after_sharding = math.ceil(len(train_dataloader) / accelerator.num_processes)
num_update_steps_per_epoch = math.ceil(len_train_dataloader_after_sharding / args.gradient_accumulation_steps)
num_training_steps_for_scheduler = (
args.num_train_epochs * num_update_steps_per_epoch * accelerator.num_processes
)
else:
num_training_steps_for_scheduler = args.max_train_steps * accelerator.num_processes
lr_scheduler = get_scheduler(
args.lr_scheduler,
optimizer=optimizer,
num_warmup_steps=args.lr_warmup_steps * accelerator.num_processes,
num_training_steps=args.max_train_steps * accelerator.num_processes,
num_warmup_steps=num_warmup_steps_for_scheduler,
num_training_steps=num_training_steps_for_scheduler,
)
# Prepare everything with our `accelerator`.
@@ -866,8 +871,14 @@ def main():
# We need to recalculate our total training steps as the size of the training dataloader may have changed.
num_update_steps_per_epoch = math.ceil(len(train_dataloader) / args.gradient_accumulation_steps)
if overrode_max_train_steps:
if args.max_train_steps is None:
args.max_train_steps = args.num_train_epochs * num_update_steps_per_epoch
if num_training_steps_for_scheduler != args.max_train_steps * accelerator.num_processes:
logger.warning(
f"The length of the 'train_dataloader' after 'accelerator.prepare' ({len(train_dataloader)}) does not match "
f"the expected length ({len_train_dataloader_after_sharding}) when the learning rate scheduler was created. "
f"This inconsistency may result in the learning rate scheduler not functioning properly."
)
# Afterwards we recalculate our number of training epochs
args.num_train_epochs = math.ceil(args.max_train_steps / num_update_steps_per_epoch)

View File

@@ -478,7 +478,7 @@ def parse_args(input_args=None):
parser.add_argument(
"--debug_loss",
action="store_true",
help="debug loss for each image, if filenames are awailable in the dataset",
help="debug loss for each image, if filenames are available in the dataset",
)
if input_args is not None:

View File

@@ -23,4 +23,25 @@ accelerate launch textual_inversion_sdxl.py \
--output_dir="./textual_inversion_cat_sdxl"
```
For now, only training of the first text encoder is supported.
Training of both text encoders is supported.
### Inference Example
Once you have trained a model using above command, the inference can be done simply using the `StableDiffusionXLPipeline`.
Make sure to include the `placeholder_token` in your prompt.
```python
from diffusers import StableDiffusionXLPipeline
import torch
model_id = "./textual_inversion_cat_sdxl"
pipe = StableDiffusionXLPipeline.from_pretrained(model_id,torch_dtype=torch.float16).to("cuda")
prompt = "A <cat-toy> backpack"
image = pipe(prompt, num_inference_steps=50, guidance_scale=7.5).images[0]
image.save("cat-backpack.png")
image = pipe(prompt="", prompt_2=prompt, num_inference_steps=50, guidance_scale=7.5).images[0]
image.save("cat-backpack-prompt_2.png")
```

View File

@@ -135,7 +135,7 @@ def log_validation(
pipeline = DiffusionPipeline.from_pretrained(
args.pretrained_model_name_or_path,
text_encoder=accelerator.unwrap_model(text_encoder_1),
text_encoder_2=text_encoder_2,
text_encoder_2=accelerator.unwrap_model(text_encoder_2),
tokenizer=tokenizer_1,
tokenizer_2=tokenizer_2,
unet=unet,
@@ -678,36 +678,54 @@ def main():
f"The tokenizer already contains the token {args.placeholder_token}. Please pass a different"
" `placeholder_token` that is not already in the tokenizer."
)
num_added_tokens = tokenizer_2.add_tokens(placeholder_tokens)
if num_added_tokens != args.num_vectors:
raise ValueError(
f"The 2nd tokenizer already contains the token {args.placeholder_token}. Please pass a different"
" `placeholder_token` that is not already in the tokenizer."
)
# Convert the initializer_token, placeholder_token to ids
token_ids = tokenizer_1.encode(args.initializer_token, add_special_tokens=False)
token_ids_2 = tokenizer_2.encode(args.initializer_token, add_special_tokens=False)
# Check if initializer_token is a single token or a sequence of tokens
if len(token_ids) > 1:
if len(token_ids) > 1 or len(token_ids_2) > 1:
raise ValueError("The initializer token must be a single token.")
initializer_token_id = token_ids[0]
placeholder_token_ids = tokenizer_1.convert_tokens_to_ids(placeholder_tokens)
initializer_token_id_2 = token_ids_2[0]
placeholder_token_ids_2 = tokenizer_2.convert_tokens_to_ids(placeholder_tokens)
# Resize the token embeddings as we are adding new special tokens to the tokenizer
text_encoder_1.resize_token_embeddings(len(tokenizer_1))
text_encoder_2.resize_token_embeddings(len(tokenizer_2))
# Initialise the newly added placeholder token with the embeddings of the initializer token
token_embeds = text_encoder_1.get_input_embeddings().weight.data
token_embeds_2 = text_encoder_2.get_input_embeddings().weight.data
with torch.no_grad():
for token_id in placeholder_token_ids:
token_embeds[token_id] = token_embeds[initializer_token_id].clone()
for token_id in placeholder_token_ids_2:
token_embeds_2[token_id] = token_embeds_2[initializer_token_id_2].clone()
# Freeze vae and unet
vae.requires_grad_(False)
unet.requires_grad_(False)
text_encoder_2.requires_grad_(False)
# Freeze all parameters except for the token embeddings in text encoder
text_encoder_1.text_model.encoder.requires_grad_(False)
text_encoder_1.text_model.final_layer_norm.requires_grad_(False)
text_encoder_1.text_model.embeddings.position_embedding.requires_grad_(False)
text_encoder_2.text_model.encoder.requires_grad_(False)
text_encoder_2.text_model.final_layer_norm.requires_grad_(False)
text_encoder_2.text_model.embeddings.position_embedding.requires_grad_(False)
if args.gradient_checkpointing:
text_encoder_1.gradient_checkpointing_enable()
text_encoder_2.gradient_checkpointing_enable()
if args.enable_xformers_memory_efficient_attention:
if is_xformers_available():
@@ -746,7 +764,11 @@ def main():
optimizer_class = torch.optim.AdamW
optimizer = optimizer_class(
text_encoder_1.get_input_embeddings().parameters(), # only optimize the embeddings
# only optimize the embeddings
[
text_encoder_1.text_model.embeddings.token_embedding.weight,
text_encoder_2.text_model.embeddings.token_embedding.weight,
],
lr=args.learning_rate,
betas=(args.adam_beta1, args.adam_beta2),
weight_decay=args.adam_weight_decay,
@@ -786,9 +808,10 @@ def main():
)
text_encoder_1.train()
text_encoder_2.train()
# Prepare everything with our `accelerator`.
text_encoder_1, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
text_encoder_1, optimizer, train_dataloader, lr_scheduler
text_encoder_1, text_encoder_2, optimizer, train_dataloader, lr_scheduler = accelerator.prepare(
text_encoder_1, text_encoder_2, optimizer, train_dataloader, lr_scheduler
)
# For mixed precision training we cast all non-trainable weigths (vae, non-lora text_encoder and non-lora unet) to half-precision
@@ -866,11 +889,13 @@ def main():
# keep original embeddings as reference
orig_embeds_params = accelerator.unwrap_model(text_encoder_1).get_input_embeddings().weight.data.clone()
orig_embeds_params_2 = accelerator.unwrap_model(text_encoder_2).get_input_embeddings().weight.data.clone()
for epoch in range(first_epoch, args.num_train_epochs):
text_encoder_1.train()
text_encoder_2.train()
for step, batch in enumerate(train_dataloader):
with accelerator.accumulate(text_encoder_1):
with accelerator.accumulate([text_encoder_1, text_encoder_2]):
# Convert images to latent space
latents = vae.encode(batch["pixel_values"].to(dtype=weight_dtype)).latent_dist.sample().detach()
latents = latents * vae.config.scaling_factor
@@ -892,9 +917,7 @@ def main():
.hidden_states[-2]
.to(dtype=weight_dtype)
)
encoder_output_2 = text_encoder_2(
batch["input_ids_2"].reshape(batch["input_ids_1"].shape[0], -1), output_hidden_states=True
)
encoder_output_2 = text_encoder_2(batch["input_ids_2"], output_hidden_states=True)
encoder_hidden_states_2 = encoder_output_2.hidden_states[-2].to(dtype=weight_dtype)
original_size = [
(batch["original_size"][0][i].item(), batch["original_size"][1][i].item())
@@ -938,11 +961,16 @@ def main():
# Let's make sure we don't update any embedding weights besides the newly added token
index_no_updates = torch.ones((len(tokenizer_1),), dtype=torch.bool)
index_no_updates[min(placeholder_token_ids) : max(placeholder_token_ids) + 1] = False
index_no_updates_2 = torch.ones((len(tokenizer_2),), dtype=torch.bool)
index_no_updates_2[min(placeholder_token_ids_2) : max(placeholder_token_ids_2) + 1] = False
with torch.no_grad():
accelerator.unwrap_model(text_encoder_1).get_input_embeddings().weight[
index_no_updates
] = orig_embeds_params[index_no_updates]
accelerator.unwrap_model(text_encoder_2).get_input_embeddings().weight[
index_no_updates_2
] = orig_embeds_params_2[index_no_updates_2]
# Checks if the accelerator has performed an optimization step behind the scenes
if accelerator.sync_gradients:
@@ -960,6 +988,16 @@ def main():
save_path,
safe_serialization=True,
)
weight_name = f"learned_embeds_2-steps-{global_step}.safetensors"
save_path = os.path.join(args.output_dir, weight_name)
save_progress(
text_encoder_2,
placeholder_token_ids_2,
accelerator,
args,
save_path,
safe_serialization=True,
)
if accelerator.is_main_process:
if global_step % args.checkpointing_steps == 0:
@@ -1034,7 +1072,7 @@ def main():
pipeline = DiffusionPipeline.from_pretrained(
args.pretrained_model_name_or_path,
text_encoder=accelerator.unwrap_model(text_encoder_1),
text_encoder_2=text_encoder_2,
text_encoder_2=accelerator.unwrap_model(text_encoder_2),
vae=vae,
unet=unet,
tokenizer=tokenizer_1,
@@ -1052,6 +1090,16 @@ def main():
save_path,
safe_serialization=True,
)
weight_name = "learned_embeds_2.safetensors"
save_path = os.path.join(args.output_dir, weight_name)
save_progress(
text_encoder_2,
placeholder_token_ids_2,
accelerator,
args,
save_path,
safe_serialization=True,
)
if args.push_to_hub:
save_model_card(

View File

@@ -263,6 +263,41 @@ class ModelMixin(torch.nn.Module, PushToHubMixin):
"""
self.set_use_memory_efficient_attention_xformers(False)
def enable_dynamic_upcasting(self, upcast_dtype=None):
upcast_dtype = upcast_dtype or torch.float32
downcast_dtype = self.dtype
def upcast_hook_fn(module):
module = module.to(upcast_dtype)
def downcast_hook_fn(module):
module = module.to(downcast_dtype)
def fn_recursive_upcast(module):
has_children = list(module.children())
if not has_children:
module.register_forward_pre_hook(upcast_hook_fn)
module.register_forward_hook(downcast_hook_fn)
for child in module.children():
fn_recursive_upcast(child)
for module in self.children():
fn_recursive_upcast(module)
def disable_dynamic_upcasting(self):
def fn_recursive_upcast(module):
has_children = list(module.children())
if not has_children:
module._forward_pre_hooks = OrderedDict()
module._forward_hooks = OrderedDict()
for child in module.children():
fn_recursive_upcast(child)
for module in self.children():
fn_recursive_upcast(module)
def save_pretrained(
self,
save_directory: Union[str, os.PathLike],

View File

@@ -311,27 +311,6 @@ class FluxTransformer2DModel(ModelMixin, ConfigMixin, PeftAdapterMixin, FromOrig
if hasattr(module, "gradient_checkpointing"):
module.gradient_checkpointing = value
def enable_mixed_precision_inference(self, upcast_dtype=None, downcast_dtype=None):
downcast_dtype = downcast_dtype or torch.float16
upcast_dtype = upcast_dtype or torch.float32
def pre_hook_fn(module, input):
module = module.to(upcast_dtype)
def hook_fn(module, input, output):
module = module.to(downcast_dtype)
def fn_recursive_upcast(module: torch.nn.Module):
if not list(module.children()):
module.register_forward_pre_hook(pre_hook_fn)
module.register_forward_hook(hook_fn)
for child in module.children():
fn_recursive_upcast(child)
for module in self.children():
fn_recursive_upcast(module)
def forward(
self,
hidden_states: torch.Tensor,

View File

@@ -19,6 +19,7 @@ import inspect
import os
import re
import sys
from collections import OrderedDict
from dataclasses import dataclass
from pathlib import Path
from typing import Any, Callable, Dict, List, Optional, Union, get_args, get_origin
@@ -1172,6 +1173,93 @@ class DiffusionPipeline(ConfigMixin, PushToHubMixin):
component.to("cpu")
self.hf_device_map = None
def enable_dynamic_upcasting(
self,
components: Optional[List[str]] = None,
upcast_dtype: Optional[torch.dtype] = None,
):
r"""
Enable module-wise dynamic upcasting. This allows models to be loaded into the GPU in a low memory dtype e.g.
torch.float8_e4m3fn, but perform inference using a dtype that is supported on the GPU, by casting the module to
the appropriate dtype right before the foward pass. The module is then moved back to the low memory dtype after
the foward pass.
"""
if components is None:
raise ValueError("Please provide a list of pipeline component names to apply dynamic upcasting")
def fn_recursive_upcast(module, dtype, original_dtype, keep_in_fp32_modules):
has_children = list(module.children())
upcast_dtype = dtype
downcast_dtype = original_dtype
def upcast_hook_fn(module, inputs):
module = module.to(upcast_dtype)
def downcast_hook_fn(module, *args, **kwargs):
module = module.to(downcast_dtype)
if not has_children:
module.register_forward_pre_hook(upcast_hook_fn)
module.register_forward_hook(downcast_hook_fn)
for name, child in module.named_children():
if any(module_to_keep_in_fp32 in name.split(".") for module_to_keep_in_fp32 in keep_in_fp32_modules):
dtype = torch.float32
else:
dtype = upcast_dtype
fn_recursive_upcast(child, dtype, original_dtype, keep_in_fp32_modules)
for component in components:
if not hasattr(self, component):
raise ValueError(f"Pipeline has no component named: {component}")
component_module = getattr(self, component)
if not isinstance(component_module, torch.nn.Module):
raise ValueError(
f"Pipeline component: {component} is not a torch.nn.Module. Cannot apply dynamic upcasting."
)
use_keep_in_fp32_modules = (
hasattr(component_module, "_keep_in_fp32_modules")
and (component_module._keep_in_fp32_modules is not None)
and (upcast_dtype != torch.float32)
)
if use_keep_in_fp32_modules:
keep_in_fp32_modules = component_module._keep_in_fp32_modules
else:
keep_in_fp32_modules = []
original_dtype = component_module.dtype
for name, module in component_module.named_children():
fn_recursive_upcast(module, upcast_dtype, original_dtype, keep_in_fp32_modules)
def disable_dynamic_upcasting(
self,
):
def fn_recursive_upcast(module):
has_children = list(module.children())
if not has_children:
module._forward_pre_hooks = OrderedDict()
module._forward_hooks = OrderedDict()
for child in module.children():
fn_recursive_upcast(child)
for component in self.components:
if not hasattr(self, component):
raise ValueError(f"Pipeline has no component named: {component}")
component_module = getattr(self, component)
if not issubclass(component_module, torch.nn.Module):
raise ValueError(
f"Pipeline component: {component} is not an torch.nn.Module. Cannot apply dynamic upcasting."
)
for module in component_module.children():
fn_recursive_upcast(module)
@classmethod
@validate_hf_hub_args
def download(cls, pretrained_model_name, **kwargs) -> Union[str, os.PathLike]:

View File

@@ -9,7 +9,7 @@ import numpy as np
import PIL.Image
import PIL.ImageOps
from .import_utils import BACKENDS_MAPPING, is_opencv_available
from .import_utils import BACKENDS_MAPPING, is_imageio_available, is_opencv_available
from .logging import get_logger
@@ -112,9 +112,9 @@ def export_to_obj(mesh, output_obj_path: str = None):
f.writelines("\n".join(combined_data))
def export_to_video(
def _legacy_export_to_video(
video_frames: Union[List[np.ndarray], List[PIL.Image.Image]], output_video_path: str = None, fps: int = 10
) -> str:
):
if is_opencv_available():
import cv2
else:
@@ -134,4 +134,51 @@ def export_to_video(
for i in range(len(video_frames)):
img = cv2.cvtColor(video_frames[i], cv2.COLOR_RGB2BGR)
video_writer.write(img)
return output_video_path
def export_to_video(
video_frames: Union[List[np.ndarray], List[PIL.Image.Image]], output_video_path: str = None, fps: int = 10
) -> str:
# TODO: Dhruv. Remove by Diffusers release 0.33.0
# Added to prevent breaking existing code
if not is_imageio_available():
logger.warning(
(
"It is recommended to use `export_to_video` with `imageio` and `imageio-ffmpeg` as a backend. \n"
"These libraries are not present in your environment. Attempting to use legacy OpenCV backend to export video. \n"
"Support for the OpenCV backend will be deprecated in a future Diffusers version"
)
)
return _legacy_export_to_video(video_frames, output_video_path, fps)
if is_imageio_available():
import imageio
else:
raise ImportError(BACKENDS_MAPPING["imageio"][1].format("export_to_video"))
try:
imageio.plugins.ffmpeg.get_exe()
except AttributeError:
raise AttributeError(
(
"Found an existing imageio backend in your environment. Attempting to export video with imageio. \n"
"Unable to find a compatible ffmpeg installation in your environment to use with imageio. Please install via `pip install imageio-ffmpeg"
)
)
if output_video_path is None:
output_video_path = tempfile.NamedTemporaryFile(suffix=".mp4").name
if isinstance(video_frames[0], np.ndarray):
video_frames = [(frame * 255).astype(np.uint8) for frame in video_frames]
elif isinstance(video_frames[0], PIL.Image.Image):
video_frames = [np.array(frame) for frame in video_frames]
with imageio.get_writer(output_video_path, fps=fps) as writer:
for frame in video_frames:
writer.append_data(frame)
return output_video_path

View File

@@ -330,6 +330,15 @@ except importlib_metadata.PackageNotFoundError:
_is_google_colab = "google.colab" in sys.modules or any(k.startswith("COLAB_") for k in os.environ)
_imageio_available = importlib.util.find_spec("imageio") is not None
if _imageio_available:
try:
_imageio_version = importlib_metadata.version("imageio")
logger.debug(f"Successfully imported imageio version {_imageio_version}")
except importlib_metadata.PackageNotFoundError:
_imageio_available = False
def is_torch_available():
return _torch_available
@@ -447,6 +456,10 @@ def is_sentencepiece_available():
return _sentencepiece_available
def is_imageio_available():
return _imageio_available
# docstyle-ignore
FLAX_IMPORT_ERROR = """
{0} requires the FLAX library but it was not found in your environment. Checkout the instructions on the
@@ -575,6 +588,11 @@ BITSANDBYTES_IMPORT_ERROR = """
{0} requires the bitsandbytes library but it was not found in your environment. You can install it with pip: `pip install bitsandbytes`
"""
# docstyle-ignore
IMAGEIO_IMPORT_ERROR = """
{0} requires the imageio library and ffmpeg but it was not found in your environment. You can install it with pip: `pip install imageio imageio-ffmpeg`
"""
BACKENDS_MAPPING = OrderedDict(
[
("bs4", (is_bs4_available, BS4_IMPORT_ERROR)),
@@ -599,6 +617,7 @@ BACKENDS_MAPPING = OrderedDict(
("safetensors", (is_safetensors_available, SAFETENSORS_IMPORT_ERROR)),
("bitsandbytes", (is_bitsandbytes_available, BITSANDBYTES_IMPORT_ERROR)),
("sentencepiece", (is_sentencepiece_available, SENTENCEPIECE_IMPORT_ERROR)),
("imageio", (is_imageio_available, IMAGEIO_IMPORT_ERROR)),
]
)

View File

@@ -6,7 +6,7 @@ import PIL.Image
import PIL.ImageOps
import requests
from .import_utils import BACKENDS_MAPPING, is_opencv_available
from .import_utils import BACKENDS_MAPPING, is_imageio_available
def load_image(
@@ -81,7 +81,8 @@ def load_video(
if is_url:
video_data = requests.get(video, stream=True).raw
video_path = tempfile.NamedTemporaryFile(suffix=os.path.splitext(video)[1], delete=False).name
suffix = os.path.splitext(video)[1] or ".mp4"
video_path = tempfile.NamedTemporaryFile(suffix=suffix, delete=False).name
was_tempfile_created = True
with open(video_path, "wb") as f:
f.write(video_data.read())
@@ -99,19 +100,22 @@ def load_video(
pass
else:
if is_opencv_available():
import cv2
if is_imageio_available():
import imageio
else:
raise ImportError(BACKENDS_MAPPING["opencv"][1].format("load_video"))
raise ImportError(BACKENDS_MAPPING["imageio"][1].format("load_video"))
video_capture = cv2.VideoCapture(video)
success, frame = video_capture.read()
while success:
frame = cv2.cvtColor(frame, cv2.COLOR_BGR2RGB)
pil_images.append(PIL.Image.fromarray(frame))
success, frame = video_capture.read()
try:
imageio.plugins.ffmpeg.get_exe()
except AttributeError:
raise AttributeError(
"`Unable to find an ffmpeg installation on your machine. Please install via `pip install imageio-ffmpeg"
)
video_capture.release()
with imageio.get_reader(video) as reader:
# Read all frames
for frame in reader:
pil_images.append(PIL.Image.fromarray(frame))
if was_tempfile_created:
os.remove(video_path)