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47 Commits

Author SHA1 Message Date
Dhruv Nair
390d0879fe update 2024-07-15 12:57:28 +00:00
Dhruv Nair
1fb8082af9 Merge branch 'main' into shm-size 2024-07-15 12:51:14 +00:00
Dhruv Nair
54e6f2d99e Merge branch 'main' into shm-size 2024-07-04 09:12:16 +00:00
Dhruv Nair
22326b4971 update 2024-06-04 08:33:22 +00:00
Dhruv Nair
7cc0ba0070 update 2024-06-03 04:52:47 +00:00
Dhruv Nair
e4bd3f5a05 update 2024-06-03 04:51:32 +00:00
Dhruv Nair
25d0277911 change runner type 2024-05-07 07:21:49 +00:00
Dhruv Nair
14d7fe3f9e update 2024-05-06 17:21:06 +00:00
Dhruv Nair
2f4e29d179 update 2024-05-06 16:28:50 +00:00
Dhruv Nair
b28311d9e5 update 2024-05-06 16:12:36 +00:00
Dhruv Nair
8a1020c91e update 2024-05-06 16:07:35 +00:00
Dhruv Nair
2483c57c44 update 2024-05-06 15:21:02 +00:00
Dhruv Nair
f18942b3e6 update 2024-05-06 14:24:29 +00:00
Dhruv Nair
7021a35eeb update 2024-05-06 13:06:44 +00:00
Dhruv Nair
e0455537c2 update 2024-05-06 09:00:56 +00:00
Dhruv Nair
c61360099b update 2024-05-06 08:36:40 +00:00
Dhruv Nair
b585832d4e update 2024-05-06 08:19:11 +00:00
Dhruv Nair
963c73c229 update 2024-05-06 06:10:01 +00:00
Dhruv Nair
09c595f212 update 2024-05-06 06:06:27 +00:00
Dhruv Nair
b1c5030418 update 2024-05-06 06:01:03 +00:00
Dhruv Nair
17e801ccb9 Merge branch 'main' into shm-size 2024-05-06 05:59:20 +00:00
Dhruv Nair
468008f6b7 update 2024-04-30 12:45:45 +00:00
Dhruv Nair
10e7d39f58 update 2024-04-30 12:40:00 +00:00
Dhruv Nair
49c01d4a4b update 2024-04-30 12:37:45 +00:00
Dhruv Nair
f2e6c24df2 update 2024-04-30 12:35:56 +00:00
Dhruv Nair
ce659bc586 Merge branch 'main' into shm-size 2024-04-01 13:03:06 +00:00
Dhruv Nair
615c12ab68 Merge branch 'main' into shm-size 2024-03-28 05:15:41 +00:00
Dhruv Nair
e34d9f1949 update 2024-03-22 13:37:48 +00:00
Dhruv Nair
4106e3f182 update 2024-03-22 12:45:44 +00:00
Dhruv Nair
1651c9a1a1 update 2024-03-22 12:08:44 +00:00
Dhruv Nair
4bfdb34b36 update 2024-03-22 11:39:09 +00:00
Dhruv Nair
3d4f987cc9 update 2024-03-22 11:34:21 +00:00
Dhruv Nair
1feac0469b update 2024-03-22 11:25:20 +00:00
Dhruv Nair
75e5cd046b update 2024-03-22 09:32:23 +00:00
Dhruv Nair
865b6638f7 update 2024-03-22 09:26:12 +00:00
Dhruv Nair
428c952289 update 2024-03-22 08:11:16 +00:00
Dhruv Nair
d16c921346 update 2024-03-22 08:10:50 +00:00
Dhruv Nair
c46380165a update 2024-03-22 07:42:00 +00:00
Dhruv Nair
a6a89aa199 update 2024-03-22 07:26:42 +00:00
Dhruv Nair
9e9a49ca24 update 2024-03-21 08:39:27 +00:00
Dhruv Nair
b9c90f7e22 update 2024-03-21 08:24:15 +00:00
Dhruv Nair
19fc3281c5 update 2024-03-21 08:18:44 +00:00
Dhruv Nair
fe8f0c9a76 update 2024-03-21 07:08:14 +00:00
Dhruv Nair
dc0bef11bd Merge branch 'shm-size' of https://github.com/huggingface/diffusers into shm-size 2024-03-20 11:09:30 +00:00
Dhruv Nair
27ef43f043 update 2024-03-20 11:09:21 +00:00
Dhruv Nair
5d5c61bd09 Merge branch 'main' into shm-size 2024-03-20 16:36:40 +05:30
Dhruv Nair
7c38bff847 update 2024-03-20 10:36:10 +00:00
939 changed files with 13773 additions and 161829 deletions

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@@ -73,7 +73,7 @@ body:
- ControlNet @sayakpaul @yiyixuxu @DN6
- T2I Adapter @sayakpaul @yiyixuxu @DN6
- IF @DN6
- Text-to-Video / Video-to-Video @DN6 @a-r-r-o-w
- Text-to-Video / Video-to-Video @DN6 @sayakpaul
- Wuerstchen @DN6
- Other: @yiyixuxu @DN6
- Improving generation quality: @asomoza

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@@ -49,7 +49,6 @@ Core library:
Integrations:
- deepspeed: HF Trainer/Accelerate: @SunMarc
- PEFT: @sayakpaul @BenjaminBossan
HF projects:

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@@ -7,24 +7,22 @@ on:
env:
DIFFUSERS_IS_CI: yes
HF_HUB_ENABLE_HF_TRANSFER: 1
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
jobs:
torch_pipelines_cuda_benchmark_tests:
env:
env:
SLACK_WEBHOOK_URL: ${{ secrets.SLACK_WEBHOOK_URL_BENCHMARK }}
name: Torch Core Pipelines CUDA Benchmarking Tests
strategy:
fail-fast: false
max-parallel: 1
runs-on:
group: aws-g6-4xlarge-plus
runs-on: [single-gpu, nvidia-gpu, a10, ci]
container:
image: diffusers/diffusers-pytorch-compile-cuda
options: --shm-size "16gb" --ipc host --gpus 0
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -51,7 +49,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: benchmark_test_reports
path: benchmarks/benchmark_outputs
@@ -60,7 +58,7 @@ jobs:
if: ${{ success() }}
run: |
pip install requests && python utils/notify_benchmarking_status.py --status=success
- name: Report failure status
if: ${{ failure() }}
run: |

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@@ -20,8 +20,7 @@ env:
jobs:
test-build-docker-images:
runs-on:
group: aws-general-8-plus
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
if: github.event_name == 'pull_request'
steps:
- name: Set up Docker Buildx
@@ -51,8 +50,7 @@ jobs:
if: steps.file_changes.outputs.all != ''
build-and-push-docker-images:
runs-on:
group: aws-general-8-plus
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
if: github.event_name != 'pull_request'
permissions:
@@ -100,4 +98,4 @@ jobs:
slack_channel: ${{ env.CI_SLACK_CHANNEL }}
title: "🤗 Results of the ${{ matrix.image-name }} Docker Image build"
status: ${{ job.status }}
slack_token: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
slack_token: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}

View File

@@ -24,8 +24,8 @@ jobs:
mirror_community_pipeline:
env:
SLACK_WEBHOOK_URL: ${{ secrets.SLACK_WEBHOOK_URL_COMMUNITY_MIRROR }}
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
steps:
# Checkout to correct ref
# If workflow dispatch
@@ -95,7 +95,7 @@ jobs:
if: ${{ success() }}
run: |
pip install requests && python utils/notify_community_pipelines_mirror.py --status=success
- name: Report failure status
if: ${{ failure() }}
run: |

View File

@@ -7,7 +7,7 @@ on:
env:
DIFFUSERS_IS_CI: yes
HF_HUB_ENABLE_HF_TRANSFER: 1
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
@@ -18,11 +18,8 @@ env:
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
name: Setup Torch Pipelines Matrix
runs-on: diffusers/diffusers-pytorch-cpu
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
@@ -30,9 +27,13 @@ jobs:
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
pip install -e .[test]
pip install -e .
pip install huggingface_hub
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
@@ -43,24 +44,22 @@ jobs:
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: test-pipelines.json
path: reports
run_nightly_tests_for_torch_pipelines:
name: Nightly Torch Pipelines CUDA Tests
name: Torch Pipelines CUDA Nightly Tests
needs: setup_torch_cuda_pipeline_matrix
strategy:
fail-fast: false
max-parallel: 8
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on:
group: aws-g4dn-2xlarge
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -68,18 +67,21 @@ jobs:
fetch-depth: 2
- name: NVIDIA-SMI
run: nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: |
python utils/print_env.py
- name: Pipeline CUDA Test
- name: Nightly PyTorch CUDA checkpoint (pipelines) tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
@@ -88,38 +90,38 @@ jobs:
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
--report-log=tests_pipeline_${{ matrix.module }}_cuda.log \
tests/pipelines/${{ matrix.module }}
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: pipeline_${{ matrix.module }}_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_tests_for_other_torch_modules:
name: Nightly Torch CUDA Tests
runs-on:
group: aws-g4dn-2xlarge
name: Torch Non-Pipelines CUDA Nightly Tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
strategy:
fail-fast: false
max-parallel: 2
matrix:
module: [models, schedulers, lora, others, single_file, examples]
module: [models, schedulers, others, examples]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -130,16 +132,16 @@ jobs:
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly PyTorch CUDA tests for non-pipeline modules
if: ${{ matrix.module != 'examples'}}
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
@@ -152,10 +154,11 @@ jobs:
- name: Run nightly example tests with Torch
if: ${{ matrix.module == 'examples' }}
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v --make-reports=examples_torch_cuda \
--report-log=examples_torch_cuda.log \
@@ -169,7 +172,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: torch_${{ matrix.module }}_cuda_test_reports
path: reports
@@ -178,73 +181,14 @@ jobs:
if: always()
run: |
pip install slack_sdk tabulate
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_big_gpu_torch_tests:
name: Torch tests on big GPU
strategy:
fail-fast: false
max-parallel: 2
runs-on:
group: aws-g6e-xlarge-plus
run_lora_nightly_tests:
name: Nightly LoRA Tests with PEFT and TORCH
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: |
python utils/print_env.py
- name: Selected Torch CUDA Test on big GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
BIG_GPU_MEMORY: 40
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-m "big_gpu_with_torch_cuda" \
--make-reports=tests_big_gpu_torch_cuda \
--report-log=tests_big_gpu_torch_cuda.log \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_big_gpu_torch_cuda_stats.txt
cat reports/tests_big_gpu_torch_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: torch_cuda_big_gpu_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
run_flax_tpu_tests:
name: Nightly Flax TPU Tests
runs-on:
group: gcp-ct5lp-hightpu-8t
if: github.event_name == 'schedule'
container:
image: diffusers/diffusers-flax-tpu
options: --shm-size "16gb" --ipc host --privileged ${{ vars.V5_LITEPOD_8_ENV}} -v /mnt/hf_cache:/mnt/hf_cache
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
@@ -258,7 +202,66 @@ jobs:
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly LoRA tests with PEFT and Torch
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_lora_cuda \
--report-log=tests_torch_lora_cuda.log \
tests/lora
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_lora_cuda_stats.txt
cat reports/tests_torch_lora_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_lora_cuda_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_flax_tpu_tests:
name: Nightly Flax TPU Tests
runs-on: docker-tpu
if: github.event_name == 'schedule'
container:
image: diffusers/diffusers-flax-tpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --privileged
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
@@ -266,7 +269,7 @@ jobs:
- name: Run nightly Flax TPU tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
@@ -282,7 +285,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: flax_tpu_test_reports
path: reports
@@ -291,15 +294,14 @@ jobs:
if: always()
run: |
pip install slack_sdk tabulate
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_onnx_tests:
name: Nightly ONNXRuntime CUDA tests on Ubuntu
runs-on:
group: aws-g4dn-2xlarge
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-onnxruntime-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
steps:
- name: Checkout diffusers
@@ -314,14 +316,15 @@ jobs:
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run Nightly ONNXRuntime CUDA tests
- name: Run nightly ONNXRuntime CUDA tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
@@ -337,189 +340,75 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: tests_onnx_cuda_reports
name: ${{ matrix.config.report }}_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_tests_apple_m1:
name: Nightly PyTorch MPS tests on MacOS
runs-on: [ self-hosted, apple-m1 ]
if: github.event_name == 'schedule'
run_nightly_quantization_tests:
name: Torch quantization nightly tests
strategy:
fail-fast: false
max-parallel: 2
matrix:
config:
- backend: "bitsandbytes"
test_location: "bnb"
- backend: "gguf"
test_location: "gguf"
runs-on:
group: aws-g6e-xlarge-plus
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "20gb" --ipc host --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: nvidia-smi
- name: Clean checkout
shell: arch -arch arm64 bash {0}
run: |
git clean -fxd
- name: Setup miniconda
uses: ./.github/actions/setup-miniconda
with:
python-version: 3.9
- name: Install dependencies
shell: arch -arch arm64 bash {0}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install -U ${{ matrix.config.backend }}
python -m uv pip install pytest-reportlog
${CONDA_RUN} python -m pip install --upgrade pip uv
${CONDA_RUN} python -m uv pip install -e [quality,test]
${CONDA_RUN} python -m uv pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
${CONDA_RUN} python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate
${CONDA_RUN} python -m uv pip install pytest-reportlog
- name: Environment
shell: arch -arch arm64 bash {0}
run: |
python utils/print_env.py
- name: ${{ matrix.config.backend }} quantization tests on GPU
${CONDA_RUN} python utils/print_env.py
- name: Run nightly PyTorch tests on M1 (MPS)
shell: arch -arch arm64 bash {0}
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
BIG_GPU_MEMORY: 40
HF_HOME: /System/Volumes/Data/mnt/cache
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_${{ matrix.config.backend }}_torch_cuda \
--report-log=tests_${{ matrix.config.backend }}_torch_cuda.log \
tests/quantization/${{ matrix.config.test_location }}
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
--report-log=tests_torch_mps.log \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_${{ matrix.config.backend }}_torch_cuda_stats.txt
cat reports/tests_${{ matrix.config.backend }}_torch_cuda_failures_short.txt
run: cat reports/tests_torch_mps_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: torch_cuda_${{ matrix.config.backend }}_reports
name: torch_mps_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
# M1 runner currently not well supported
# TODO: (Dhruv) add these back when we setup better testing for Apple Silicon
# run_nightly_tests_apple_m1:
# name: Nightly PyTorch MPS tests on MacOS
# runs-on: [ self-hosted, apple-m1 ]
# if: github.event_name == 'schedule'
#
# steps:
# - name: Checkout diffusers
# uses: actions/checkout@v3
# with:
# fetch-depth: 2
#
# - name: Clean checkout
# shell: arch -arch arm64 bash {0}
# run: |
# git clean -fxd
# - name: Setup miniconda
# uses: ./.github/actions/setup-miniconda
# with:
# python-version: 3.9
#
# - name: Install dependencies
# shell: arch -arch arm64 bash {0}
# run: |
# ${CONDA_RUN} python -m pip install --upgrade pip uv
# ${CONDA_RUN} python -m uv pip install -e [quality,test]
# ${CONDA_RUN} python -m uv pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
# ${CONDA_RUN} python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate
# ${CONDA_RUN} python -m uv pip install pytest-reportlog
# - name: Environment
# shell: arch -arch arm64 bash {0}
# run: |
# ${CONDA_RUN} python utils/print_env.py
# - name: Run nightly PyTorch tests on M1 (MPS)
# shell: arch -arch arm64 bash {0}
# env:
# HF_HOME: /System/Volumes/Data/mnt/cache
# HF_TOKEN: ${{ secrets.HF_TOKEN }}
# run: |
# ${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
# --report-log=tests_torch_mps.log \
# tests/
# - name: Failure short reports
# if: ${{ failure() }}
# run: cat reports/tests_torch_mps_failures_short.txt
#
# - name: Test suite reports artifacts
# if: ${{ always() }}
# uses: actions/upload-artifact@v4
# with:
# name: torch_mps_test_reports
# path: reports
#
# - name: Generate Report and Notify Channel
# if: always()
# run: |
# pip install slack_sdk tabulate
# python utils/log_reports.py >> $GITHUB_STEP_SUMMARY run_nightly_tests_apple_m1:
# name: Nightly PyTorch MPS tests on MacOS
# runs-on: [ self-hosted, apple-m1 ]
# if: github.event_name == 'schedule'
#
# steps:
# - name: Checkout diffusers
# uses: actions/checkout@v3
# with:
# fetch-depth: 2
#
# - name: Clean checkout
# shell: arch -arch arm64 bash {0}
# run: |
# git clean -fxd
# - name: Setup miniconda
# uses: ./.github/actions/setup-miniconda
# with:
# python-version: 3.9
#
# - name: Install dependencies
# shell: arch -arch arm64 bash {0}
# run: |
# ${CONDA_RUN} python -m pip install --upgrade pip uv
# ${CONDA_RUN} python -m uv pip install -e [quality,test]
# ${CONDA_RUN} python -m uv pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
# ${CONDA_RUN} python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate
# ${CONDA_RUN} python -m uv pip install pytest-reportlog
# - name: Environment
# shell: arch -arch arm64 bash {0}
# run: |
# ${CONDA_RUN} python utils/print_env.py
# - name: Run nightly PyTorch tests on M1 (MPS)
# shell: arch -arch arm64 bash {0}
# env:
# HF_HOME: /System/Volumes/Data/mnt/cache
# HF_TOKEN: ${{ secrets.HF_TOKEN }}
# run: |
# ${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
# --report-log=tests_torch_mps.log \
# tests/
# - name: Failure short reports
# if: ${{ failure() }}
# run: cat reports/tests_torch_mps_failures_short.txt
#
# - name: Test suite reports artifacts
# if: ${{ always() }}
# uses: actions/upload-artifact@v4
# with:
# name: torch_mps_test_reports
# path: reports
#
# - name: Generate Report and Notify Channel
# if: always()
# run: |
# pip install slack_sdk tabulate
# python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY

View File

@@ -7,7 +7,7 @@ on:
jobs:
build:
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3

View File

@@ -16,7 +16,7 @@ concurrency:
jobs:
check_dependencies:
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python

View File

@@ -16,7 +16,7 @@ concurrency:
jobs:
check_flax_dependencies:
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python

View File

@@ -15,8 +15,7 @@ concurrency:
jobs:
setup_pr_tests:
name: Setup PR Tests
runs-on:
group: aws-general-8-plus
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
container:
image: diffusers/diffusers-pytorch-cpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
@@ -74,8 +73,7 @@ jobs:
max-parallel: 2
matrix:
modules: ${{ fromJson(needs.setup_pr_tests.outputs.matrix) }}
runs-on:
group: aws-general-8-plus
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
container:
image: diffusers/diffusers-pytorch-cpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
@@ -125,13 +123,12 @@ jobs:
config:
- name: Hub tests for models, schedulers, and pipelines
framework: hub_tests_pytorch
runner: aws-general-8-plus
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
image: diffusers/diffusers-pytorch-cpu
report: torch_hub
name: ${{ matrix.config.name }}
runs-on:
group: ${{ matrix.config.runner }}
runs-on: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
@@ -171,7 +168,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: pr_${{ matrix.config.report }}_test_reports
path: reports

View File

@@ -0,0 +1,131 @@
name: Fast tests for PRs - PEFT backend
on:
pull_request:
branches:
- main
paths:
- "src/diffusers/**.py"
- "tests/**.py"
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
env:
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 4
MKL_NUM_THREADS: 4
PYTEST_TIMEOUT: 60
jobs:
check_code_quality:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check quality
run: make quality
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Quality check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make style && make quality'" >> $GITHUB_STEP_SUMMARY
check_repository_consistency:
needs: check_code_quality
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check repo consistency
run: |
python utils/check_copies.py
python utils/check_dummies.py
make deps_table_check_updated
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Repo consistency check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make fix-copies'" >> $GITHUB_STEP_SUMMARY
run_fast_tests:
needs: [check_code_quality, check_repository_consistency]
strategy:
fail-fast: false
matrix:
lib-versions: ["main", "latest"]
name: LoRA - ${{ matrix.lib-versions }}
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
container:
image: diffusers/diffusers-pytorch-cpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
if [ "${{ matrix.lib-versions }}" == "main" ]; then
python -m pip install -U peft@git+https://github.com/huggingface/peft.git
python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git
python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
else
python -m uv pip install -U peft transformers accelerate
fi
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run fast PyTorch LoRA CPU tests with PEFT backend
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
-s -v \
--make-reports=tests_${{ matrix.config.report }} \
tests/lora/
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
-s -v \
--make-reports=tests_models_lora_${{ matrix.config.report }} \
tests/models/ -k "lora"
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_${{ matrix.config.report }}_failures_short.txt
cat reports/tests_models_lora_${{ matrix.config.report }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: pr_${{ matrix.config.report }}_test_reports
path: reports

View File

@@ -22,14 +22,13 @@ concurrency:
env:
DIFFUSERS_IS_CI: yes
HF_HUB_ENABLE_HF_TRANSFER: 1
OMP_NUM_THREADS: 4
MKL_NUM_THREADS: 4
PYTEST_TIMEOUT: 60
jobs:
check_code_quality:
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
@@ -49,7 +48,7 @@ jobs:
check_repository_consistency:
needs: check_code_quality
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
@@ -78,29 +77,28 @@ jobs:
config:
- name: Fast PyTorch Pipeline CPU tests
framework: pytorch_pipelines
runner: aws-highmemory-32-plus
runner: [ self-hosted, intel-cpu, 32-cpu, 256-ram, ci ]
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_pipelines
- name: Fast PyTorch Models & Schedulers CPU tests
framework: pytorch_models
runner: aws-general-8-plus
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_models_schedulers
- name: Fast Flax CPU tests
framework: flax
runner: aws-general-8-plus
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
image: diffusers/diffusers-flax-cpu
report: flax_cpu
- name: PyTorch Example CPU tests
framework: pytorch_examples
runner: aws-general-8-plus
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
image: diffusers/diffusers-pytorch-cpu
report: torch_example_cpu
name: ${{ matrix.config.name }}
runs-on:
group: ${{ matrix.config.runner }}
runs-on: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
@@ -169,9 +167,9 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: pr_${{ matrix.config.framework }}_${{ matrix.config.report }}_test_reports
name: pr_${{ matrix.config.report }}_test_reports
path: reports
run_staging_tests:
@@ -182,8 +180,7 @@ jobs:
config:
- name: Hub tests for models, schedulers, and pipelines
framework: hub_tests_pytorch
runner:
group: aws-general-8-plus
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
image: diffusers/diffusers-pytorch-cpu
report: torch_hub
@@ -230,71 +227,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: pr_${{ matrix.config.report }}_test_reports
path: reports
run_lora_tests:
needs: [check_code_quality, check_repository_consistency]
strategy:
fail-fast: false
name: LoRA tests with PEFT main
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
# TODO (sayakpaul, DN6): revisit `--no-deps`
python -m pip install -U peft@git+https://github.com/huggingface/peft.git --no-deps
python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git --no-deps
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run fast PyTorch LoRA tests with PEFT
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
-s -v \
--make-reports=tests_peft_main \
tests/lora/
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
-s -v \
--make-reports=tests_models_lora_peft_main \
tests/models/ -k "lora"
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_lora_failures_short.txt
cat reports/tests_models_lora_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: pr_main_test_reports
path: reports

View File

@@ -16,7 +16,7 @@ concurrency:
jobs:
check_torch_dependencies:
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python

100
.github/workflows/push_check.yml vendored Normal file
View File

@@ -0,0 +1,100 @@
name: Slow Test Memory Checks
on:
push:
branches: [ shm-size ]
env:
DIFFUSERS_IS_CI: yes
HF_HUB_ENABLE_HF_TRANSFER: 1
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
PIPELINE_USAGE_CUTOFF: 50000
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
container:
image: diffusers/diffusers-pytorch-cpu
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
- name: Environment
run: |
python utils/print_env.py
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
run: |
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
echo $matrix
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: test-pipelines.json
path: reports
torch_pipelines_cuda_tests:
name: Torch Pipelines CUDA Slow Tests
needs: setup_torch_cuda_pipeline_matrix
strategy:
max-parallel: 4
fail-fast: false
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install hf_transfer
- name: Environment
run: |
python utils/print_env.py
- name: Slow PyTorch CUDA checkpoint tests on Ubuntu
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
tests/pipelines/${{ matrix.module }}
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: pipeline_${{ matrix.module }}_test_reports
path: reports

View File

@@ -1,7 +1,6 @@
name: Fast GPU Tests on main
name: Slow Tests on main
on:
workflow_dispatch:
push:
branches:
- main
@@ -12,17 +11,17 @@ on:
env:
DIFFUSERS_IS_CI: yes
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
HF_HUB_ENABLE_HF_TRANSFER: 1
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
PIPELINE_USAGE_CUTOFF: 50000
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on:
group: aws-general-8-plus
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
container:
image: diffusers/diffusers-pytorch-cpu
outputs:
@@ -47,24 +46,23 @@ jobs:
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: test-pipelines.json
path: reports
torch_pipelines_cuda_tests:
name: Torch Pipelines CUDA Tests
name: Torch Pipelines CUDA Slow Tests
needs: setup_torch_cuda_pipeline_matrix
strategy:
fail-fast: false
max-parallel: 8
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on:
group: aws-g4dn-2xlarge
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -77,11 +75,11 @@ jobs:
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: PyTorch CUDA checkpoint tests on Ubuntu
- name: Slow PyTorch CUDA checkpoint tests on Ubuntu
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
@@ -98,24 +96,21 @@ jobs:
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: pipeline_${{ matrix.module }}_test_reports
path: reports
torch_cuda_tests:
name: Torch CUDA Tests
runs-on:
group: aws-g4dn-2xlarge
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
strategy:
fail-fast: false
max-parallel: 2
matrix:
module: [models, schedulers, lora, others, single_file]
steps:
@@ -128,14 +123,13 @@ jobs:
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Run PyTorch CUDA tests
- name: Run slow PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
@@ -143,29 +137,28 @@ jobs:
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_cuda_${{ matrix.module }} \
--make-reports=tests_torch_cuda \
tests/${{ matrix.module }}
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_cuda_${{ matrix.module }}_stats.txt
cat reports/tests_torch_cuda_${{ matrix.module }}_failures_short.txt
cat reports/tests_torch_cuda_stats.txt
cat reports/tests_torch_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: torch_cuda_test_reports_${{ matrix.module }}
name: torch_cuda_test_reports
path: reports
flax_tpu_tests:
name: Flax TPU Tests
runs-on:
group: gcp-ct5lp-hightpu-8t
peft_cuda_tests:
name: PEFT CUDA Tests
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-flax-tpu
options: --shm-size "16gb" --ipc host --privileged ${{ vars.V5_LITEPOD_8_ENV}} -v /mnt/hf_cache:/mnt/hf_cache
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
@@ -179,13 +172,68 @@ jobs:
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m pip install -U peft@git+https://github.com/huggingface/peft.git
- name: Environment
run: |
python utils/print_env.py
- name: Run Flax TPU tests
- name: Run slow PEFT CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx and not PEFTLoRALoading" \
--make-reports=tests_peft_cuda \
tests/lora/
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "lora and not Flax and not Onnx and not PEFTLoRALoading" \
--make-reports=tests_peft_cuda_models_lora \
tests/models/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_peft_cuda_stats.txt
cat reports/tests_peft_cuda_failures_short.txt
cat reports/tests_peft_cuda_models_lora_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_peft_test_reports
path: reports
flax_tpu_tests:
name: Flax TPU Tests
runs-on: docker-tpu
container:
image: diffusers/diffusers-flax-tpu
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/ --privileged
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Run slow Flax TPU tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
@@ -202,15 +250,14 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: flax_tpu_test_reports
path: reports
onnx_cuda_tests:
name: ONNX CUDA Tests
runs-on:
group: aws-g4dn-2xlarge
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-onnxruntime-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/ --gpus 0
@@ -227,13 +274,13 @@ jobs:
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Run ONNXRuntime CUDA tests
- name: Run slow ONNXRuntime CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
@@ -250,7 +297,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: onnx_cuda_test_reports
path: reports
@@ -258,12 +305,11 @@ jobs:
run_torch_compile_tests:
name: PyTorch Compile CUDA tests
runs-on:
group: aws-g4dn-2xlarge
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-compile-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
steps:
- name: Checkout diffusers
@@ -293,7 +339,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: torch_compile_test_reports
path: reports
@@ -301,12 +347,11 @@ jobs:
run_xformers_tests:
name: PyTorch xformers CUDA tests
runs-on:
group: aws-g4dn-2xlarge
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-xformers-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
steps:
- name: Checkout diffusers
@@ -335,7 +380,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: torch_xformers_test_reports
path: reports
@@ -343,12 +388,11 @@ jobs:
run_examples_tests:
name: Examples PyTorch CUDA tests on Ubuntu
runs-on:
group: aws-g4dn-2xlarge
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: diffusers/diffusers-pytorch-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
steps:
- name: Checkout diffusers
@@ -386,7 +430,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: examples_test_reports
path: reports

View File

@@ -18,7 +18,6 @@ env:
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
HF_HUB_ENABLE_HF_TRANSFER: 1
PYTEST_TIMEOUT: 600
RUN_SLOW: no
@@ -30,29 +29,28 @@ jobs:
config:
- name: Fast PyTorch CPU tests on Ubuntu
framework: pytorch
runner: aws-general-8-plus
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu
- name: Fast Flax CPU tests on Ubuntu
framework: flax
runner: aws-general-8-plus
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
image: diffusers/diffusers-flax-cpu
report: flax_cpu
- name: Fast ONNXRuntime CPU tests on Ubuntu
framework: onnxruntime
runner: aws-general-8-plus
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
image: diffusers/diffusers-onnxruntime-cpu
report: onnx_cpu
- name: PyTorch Example CPU tests on Ubuntu
framework: pytorch_examples
runner: aws-general-8-plus
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
image: diffusers/diffusers-pytorch-cpu
report: torch_example_cpu
name: ${{ matrix.config.name }}
runs-on:
group: ${{ matrix.config.runner }}
runs-on: ${{ matrix.config.runner }}
container:
image: ${{ matrix.config.image }}
@@ -120,7 +118,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: pr_${{ matrix.config.report }}_test_reports
path: reports

View File

@@ -13,7 +13,6 @@ env:
HF_HOME: /mnt/cache
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
HF_HUB_ENABLE_HF_TRANSFER: 1
PYTEST_TIMEOUT: 600
RUN_SLOW: no
@@ -70,7 +69,7 @@ jobs:
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
uses: actions/upload-artifact@v2
with:
name: pr_torch_mps_test_reports
path: reports

View File

@@ -10,7 +10,7 @@ on:
jobs:
find-and-checkout-latest-branch:
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
outputs:
latest_branch: ${{ steps.set_latest_branch.outputs.latest_branch }}
steps:
@@ -36,7 +36,7 @@ jobs:
release:
needs: find-and-checkout-latest-branch
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
steps:
- name: Checkout Repo

View File

@@ -1,389 +0,0 @@
# Duplicate workflow to push_tests.yml that is meant to run on release/patch branches as a final check
# Creating a duplicate workflow here is simpler than adding complex path/branch parsing logic to push_tests.yml
# Needs to be updated if push_tests.yml updated
name: (Release) Fast GPU Tests on main
on:
push:
branches:
- "v*.*.*-release"
- "v*.*.*-patch"
env:
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
PYTEST_TIMEOUT: 600
PIPELINE_USAGE_CUTOFF: 50000
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
- name: Environment
run: |
python utils/print_env.py
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
run: |
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
echo $matrix
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: test-pipelines.json
path: reports
torch_pipelines_cuda_tests:
name: Torch Pipelines CUDA Tests
needs: setup_torch_cuda_pipeline_matrix
strategy:
fail-fast: false
max-parallel: 8
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Slow PyTorch CUDA checkpoint tests on Ubuntu
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
tests/pipelines/${{ matrix.module }}
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: pipeline_${{ matrix.module }}_test_reports
path: reports
torch_cuda_tests:
name: Torch CUDA Tests
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
defaults:
run:
shell: bash
strategy:
fail-fast: false
max-parallel: 2
matrix:
module: [models, schedulers, lora, others, single_file]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Run PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_${{ matrix.module }}_cuda \
tests/${{ matrix.module }}
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_torch_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: torch_cuda_${{ matrix.module }}_test_reports
path: reports
flax_tpu_tests:
name: Flax TPU Tests
runs-on: docker-tpu
container:
image: diffusers/diffusers-flax-tpu
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/ --privileged
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Run slow Flax TPU tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
--make-reports=tests_flax_tpu \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_flax_tpu_stats.txt
cat reports/tests_flax_tpu_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: flax_tpu_test_reports
path: reports
onnx_cuda_tests:
name: ONNX CUDA Tests
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-onnxruntime-cuda
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Run slow ONNXRuntime CUDA tests
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_onnx_cuda \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_onnx_cuda_stats.txt
cat reports/tests_onnx_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: onnx_cuda_test_reports
path: reports
run_torch_compile_tests:
name: PyTorch Compile CUDA tests
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-compile-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test,training]
- name: Environment
run: |
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
RUN_COMPILE: yes
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "compile" --make-reports=tests_torch_compile_cuda tests/
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_torch_compile_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: torch_compile_test_reports
path: reports
run_xformers_tests:
name: PyTorch xformers CUDA tests
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-xformers-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test,training]
- name: Environment
run: |
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "xformers" --make-reports=tests_torch_xformers_cuda tests/
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_torch_xformers_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: torch_xformers_test_reports
path: reports
run_examples_tests:
name: Examples PyTorch CUDA tests on Ubuntu
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test,training]
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install timm
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=examples_torch_cuda examples/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/examples_torch_cuda_stats.txt
cat reports/examples_torch_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: examples_test_reports
path: reports

View File

@@ -26,8 +26,7 @@ env:
jobs:
run_tests:
name: "Run a test on our runner from a PR"
runs-on:
group: aws-g4dn-2xlarge
runs-on: [single-gpu, nvidia-gpu, t4, ci]
container:
image: ${{ github.event.inputs.docker_image }}
options: --gpus 0 --privileged --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
@@ -71,4 +70,4 @@ jobs:
env:
PY_TEST: ${{ github.event.inputs.test }}
run: |
pytest "$PY_TEST"
pytest "$PY_TEST"

View File

@@ -1,40 +0,0 @@
name: SSH into PR runners
on:
workflow_dispatch:
inputs:
docker_image:
description: 'Name of the Docker image'
required: true
env:
IS_GITHUB_CI: "1"
HF_HUB_READ_TOKEN: ${{ secrets.HF_HUB_READ_TOKEN }}
HF_HOME: /mnt/cache
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
RUN_SLOW: yes
jobs:
ssh_runner:
name: "SSH"
runs-on:
group: aws-highmemory-32-plus
container:
image: ${{ github.event.inputs.docker_image }}
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --privileged
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Tailscale # In order to be able to SSH when a test fails
uses: huggingface/tailscale-action@main
with:
authkey: ${{ secrets.TAILSCALE_SSH_AUTHKEY }}
slackChannel: ${{ secrets.SLACK_CIFEEDBACK_CHANNEL }}
slackToken: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
waitForSSH: true

View File

@@ -1,16 +1,11 @@
name: SSH into GPU runners
name: SSH into runners
on:
workflow_dispatch:
inputs:
runner_type:
description: 'Type of runner to test (aws-g6-4xlarge-plus: a10, aws-g4dn-2xlarge: t4, aws-g6e-xlarge-plus: L40)'
type: choice
description: 'Type of runner to test (a10 or t4)'
required: true
options:
- aws-g6-4xlarge-plus
- aws-g4dn-2xlarge
- aws-g6e-xlarge-plus
docker_image:
description: 'Name of the Docker image'
required: true
@@ -27,8 +22,7 @@ env:
jobs:
ssh_runner:
name: "SSH"
runs-on:
group: "${{ github.event.inputs.runner_type }}"
runs-on: [single-gpu, nvidia-gpu, "${{ github.event.inputs.runner_type }}", ci]
container:
image: ${{ github.event.inputs.docker_image }}
options: --shm-size "16gb" --ipc host -v /mnt/cache/.cache/huggingface/diffusers:/mnt/cache/ --gpus 0 --privileged

View File

@@ -8,10 +8,7 @@ jobs:
close_stale_issues:
name: Close Stale Issues
if: github.repository == 'huggingface/diffusers'
runs-on: ubuntu-22.04
permissions:
issues: write
pull-requests: write
runs-on: ubuntu-latest
env:
GITHUB_TOKEN: ${{ secrets.GITHUB_TOKEN }}
steps:

View File

@@ -5,7 +5,7 @@ name: Secret Leaks
jobs:
trufflehog:
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
steps:
- name: Checkout code
uses: actions/checkout@v4

View File

@@ -5,7 +5,7 @@ on:
jobs:
build:
runs-on: ubuntu-22.04
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3

View File

@@ -57,13 +57,13 @@ Any question or comment related to the Diffusers library can be asked on the [di
- ...
Every question that is asked on the forum or on Discord actively encourages the community to publicly
share knowledge and might very well help a beginner in the future who has the same question you're
share knowledge and might very well help a beginner in the future that has the same question you're
having. Please do pose any questions you might have.
In the same spirit, you are of immense help to the community by answering such questions because this way you are publicly documenting knowledge for everybody to learn from.
**Please** keep in mind that the more effort you put into asking or answering a question, the higher
the quality of the publicly documented knowledge. In the same way, well-posed and well-answered questions create a high-quality knowledge database accessible to everybody, while badly posed questions or answers reduce the overall quality of the public knowledge database.
In short, a high quality question or answer is *precise*, *concise*, *relevant*, *easy-to-understand*, *accessible*, and *well-formatted/well-posed*. For more information, please have a look through the [How to write a good issue](#how-to-write-a-good-issue) section.
In short, a high quality question or answer is *precise*, *concise*, *relevant*, *easy-to-understand*, *accessible*, and *well-formated/well-posed*. For more information, please have a look through the [How to write a good issue](#how-to-write-a-good-issue) section.
**NOTE about channels**:
[*The forum*](https://discuss.huggingface.co/c/discussion-related-to-httpsgithubcomhuggingfacediffusers/63) is much better indexed by search engines, such as Google. Posts are ranked by popularity rather than chronologically. Hence, it's easier to look up questions and answers that we posted some time ago.
@@ -503,4 +503,4 @@ $ git push --set-upstream origin your-branch-for-syncing
### Style guide
For documentation strings, 🧨 Diffusers follows the [Google style](https://google.github.io/styleguide/pyguide.html).
For documentation strings, 🧨 Diffusers follows the [Google style](https://google.github.io/styleguide/pyguide.html).

View File

@@ -15,7 +15,7 @@ specific language governing permissions and limitations under the License.
🧨 Diffusers provides **state-of-the-art** pretrained diffusion models across multiple modalities.
Its purpose is to serve as a **modular toolbox** for both inference and training.
We aim to build a library that stands the test of time and therefore take API design very seriously.
We aim at building a library that stands the test of time and therefore take API design very seriously.
In a nutshell, Diffusers is built to be a natural extension of PyTorch. Therefore, most of our design choices are based on [PyTorch's Design Principles](https://pytorch.org/docs/stable/community/design.html#pytorch-design-philosophy). Let's go over the most important ones:
@@ -65,7 +65,7 @@ Pipelines are designed to be easy to use (therefore do not follow [*Simple over
The following design principles are followed:
- Pipelines follow the single-file policy. All pipelines can be found in individual directories under src/diffusers/pipelines. One pipeline folder corresponds to one diffusion paper/project/release. Multiple pipeline files can be gathered in one pipeline folder, as its done for [`src/diffusers/pipelines/stable-diffusion`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines/stable_diffusion). If pipelines share similar functionality, one can make use of the [# Copied from mechanism](https://github.com/huggingface/diffusers/blob/125d783076e5bd9785beb05367a2d2566843a271/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_img2img.py#L251).
- Pipelines all inherit from [`DiffusionPipeline`].
- Every pipeline consists of different model and scheduler components, that are documented in the [`model_index.json` file](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5/blob/main/model_index.json), are accessible under the same name as attributes of the pipeline and can be shared between pipelines with [`DiffusionPipeline.components`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.components) function.
- Every pipeline consists of different model and scheduler components, that are documented in the [`model_index.json` file](https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/model_index.json), are accessible under the same name as attributes of the pipeline and can be shared between pipelines with [`DiffusionPipeline.components`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.components) function.
- Every pipeline should be loadable via the [`DiffusionPipeline.from_pretrained`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.from_pretrained) function.
- Pipelines should be used **only** for inference.
- Pipelines should be very readable, self-explanatory, and easy to tweak.
@@ -107,4 +107,4 @@ The following design principles are followed:
- Every scheduler exposes the timesteps to be "looped over" via a `timesteps` attribute, which is an array of timesteps the model will be called upon.
- The `step(...)` function takes a predicted model output and the "current" sample (x_t) and returns the "previous", slightly more denoised sample (x_t-1).
- Given the complexity of diffusion schedulers, the `step` function does not expose all the complexity and can be a bit of a "black box".
- In almost all cases, novel schedulers shall be implemented in a new scheduling file.
- In almost all cases, novel schedulers shall be implemented in a new scheduling file.

View File

@@ -67,13 +67,13 @@ Please refer to the [How to use Stable Diffusion in Apple Silicon](https://huggi
## Quickstart
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 30,000+ checkpoints):
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 27.000+ checkpoints):
```python
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipeline = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipeline.to("cuda")
pipeline("An image of a squirrel in Picasso style").images[0]
```
@@ -112,9 +112,9 @@ Check out the [Quickstart](https://huggingface.co/docs/diffusers/quicktour) to l
| **Documentation** | **What can I learn?** |
|---------------------------------------------------------------------|-------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
| [Tutorial](https://huggingface.co/docs/diffusers/tutorials/tutorial_overview) | A basic crash course for learning how to use the library's most important features like using models and schedulers to build your own diffusion system, and training your own diffusion model. |
| [Loading](https://huggingface.co/docs/diffusers/using-diffusers/loading) | Guides for how to load and configure all the components (pipelines, models, and schedulers) of the library, as well as how to use different schedulers. |
| [Pipelines for inference](https://huggingface.co/docs/diffusers/using-diffusers/overview_techniques) | Guides for how to use pipelines for different inference tasks, batched generation, controlling generated outputs and randomness, and how to contribute a pipeline to the library. |
| [Optimization](https://huggingface.co/docs/diffusers/optimization/fp16) | Guides for how to optimize your diffusion model to run faster and consume less memory. |
| [Loading](https://huggingface.co/docs/diffusers/using-diffusers/loading_overview) | Guides for how to load and configure all the components (pipelines, models, and schedulers) of the library, as well as how to use different schedulers. |
| [Pipelines for inference](https://huggingface.co/docs/diffusers/using-diffusers/pipeline_overview) | Guides for how to use pipelines for different inference tasks, batched generation, controlling generated outputs and randomness, and how to contribute a pipeline to the library. |
| [Optimization](https://huggingface.co/docs/diffusers/optimization/opt_overview) | Guides for how to optimize your diffusion model to run faster and consume less memory. |
| [Training](https://huggingface.co/docs/diffusers/training/overview) | Guides for how to train a diffusion model for different tasks with different training techniques. |
## Contribution
@@ -144,7 +144,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
<tr style="border-top: 2px solid black">
<td>Text-to-Image</td>
<td><a href="https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/text2img">Stable Diffusion Text-to-Image</a></td>
<td><a href="https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5"> stable-diffusion-v1-5/stable-diffusion-v1-5 </a></td>
<td><a href="https://huggingface.co/runwayml/stable-diffusion-v1-5"> runwayml/stable-diffusion-v1-5 </a></td>
</tr>
<tr>
<td>Text-to-Image</td>
@@ -174,7 +174,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
<tr>
<td>Text-guided Image-to-Image</td>
<td><a href="https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/img2img">Stable Diffusion Image-to-Image</a></td>
<td><a href="https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5"> stable-diffusion-v1-5/stable-diffusion-v1-5 </a></td>
<td><a href="https://huggingface.co/runwayml/stable-diffusion-v1-5"> runwayml/stable-diffusion-v1-5 </a></td>
</tr>
<tr style="border-top: 2px solid black">
<td>Text-guided Image Inpainting</td>
@@ -202,7 +202,6 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
- https://github.com/microsoft/TaskMatrix
- https://github.com/invoke-ai/InvokeAI
- https://github.com/InstantID/InstantID
- https://github.com/apple/ml-stable-diffusion
- https://github.com/Sanster/lama-cleaner
- https://github.com/IDEA-Research/Grounded-Segment-Anything
@@ -210,7 +209,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
- https://github.com/deep-floyd/IF
- https://github.com/bentoml/BentoML
- https://github.com/bmaltais/kohya_ss
- +14,000 other amazing GitHub repositories 💪
- +12.000 other amazing GitHub repositories 💪
Thank you for using us ❤️.

View File

@@ -34,7 +34,7 @@ from utils import ( # noqa: E402
RESOLUTION_MAPPING = {
"Lykon/DreamShaper": (512, 512),
"runwayml/stable-diffusion-v1-5": (512, 512),
"lllyasviel/sd-controlnet-canny": (512, 512),
"diffusers/controlnet-canny-sdxl-1.0": (1024, 1024),
"TencentARC/t2iadapter_canny_sd14v1": (512, 512),
@@ -268,7 +268,7 @@ class IPAdapterTextToImageBenchmark(TextToImageBenchmark):
class ControlNetBenchmark(TextToImageBenchmark):
pipeline_class = StableDiffusionControlNetPipeline
aux_network_class = ControlNetModel
root_ckpt = "Lykon/DreamShaper"
root_ckpt = "runwayml/stable-diffusion-v1-5"
url = "https://huggingface.co/datasets/diffusers/docs-images/resolve/main/benchmarking/canny_image_condition.png"
image = load_image(url).convert("RGB")
@@ -311,7 +311,7 @@ class ControlNetSDXLBenchmark(ControlNetBenchmark):
class T2IAdapterBenchmark(ControlNetBenchmark):
pipeline_class = StableDiffusionAdapterPipeline
aux_network_class = T2IAdapter
root_ckpt = "Lykon/DreamShaper"
root_ckpt = "CompVis/stable-diffusion-v1-4"
url = "https://huggingface.co/datasets/diffusers/docs-images/resolve/main/benchmarking/canny_for_adapter.png"
image = load_image(url).convert("L")

View File

@@ -7,8 +7,7 @@ from base_classes import IPAdapterTextToImageBenchmark # noqa: E402
IP_ADAPTER_CKPTS = {
# because original SD v1.5 has been taken down.
"Lykon/DreamShaper": ("h94/IP-Adapter", "ip-adapter_sd15.bin"),
"runwayml/stable-diffusion-v1-5": ("h94/IP-Adapter", "ip-adapter_sd15.bin"),
"stabilityai/stable-diffusion-xl-base-1.0": ("h94/IP-Adapter", "ip-adapter_sdxl.bin"),
}
@@ -18,7 +17,7 @@ if __name__ == "__main__":
parser.add_argument(
"--ckpt",
type=str,
default="rstabilityai/stable-diffusion-xl-base-1.0",
default="runwayml/stable-diffusion-v1-5",
choices=list(IP_ADAPTER_CKPTS.keys()),
)
parser.add_argument("--batch_size", type=int, default=1)

View File

@@ -11,9 +11,9 @@ if __name__ == "__main__":
parser.add_argument(
"--ckpt",
type=str,
default="Lykon/DreamShaper",
default="runwayml/stable-diffusion-v1-5",
choices=[
"Lykon/DreamShaper",
"runwayml/stable-diffusion-v1-5",
"stabilityai/stable-diffusion-2-1",
"stabilityai/stable-diffusion-xl-refiner-1.0",
"stabilityai/sdxl-turbo",

View File

@@ -11,9 +11,9 @@ if __name__ == "__main__":
parser.add_argument(
"--ckpt",
type=str,
default="Lykon/DreamShaper",
default="runwayml/stable-diffusion-v1-5",
choices=[
"Lykon/DreamShaper",
"runwayml/stable-diffusion-v1-5",
"stabilityai/stable-diffusion-2-1",
"stabilityai/stable-diffusion-xl-base-1.0",
],

View File

@@ -7,7 +7,7 @@ from base_classes import TextToImageBenchmark, TurboTextToImageBenchmark # noqa
ALL_T2I_CKPTS = [
"Lykon/DreamShaper",
"runwayml/stable-diffusion-v1-5",
"segmind/SSD-1B",
"stabilityai/stable-diffusion-xl-base-1.0",
"kandinsky-community/kandinsky-2-2-decoder",
@@ -21,7 +21,7 @@ if __name__ == "__main__":
parser.add_argument(
"--ckpt",
type=str,
default="Lykon/DreamShaper",
default="runwayml/stable-diffusion-v1-5",
choices=ALL_T2I_CKPTS,
)
parser.add_argument("--batch_size", type=int, default=1)

View File

@@ -3,7 +3,7 @@ import sys
import pandas as pd
from huggingface_hub import hf_hub_download, upload_file
from huggingface_hub.utils import EntryNotFoundError
from huggingface_hub.utils._errors import EntryNotFoundError
sys.path.append(".")

View File

@@ -43,7 +43,6 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
numpy==1.26.4 \
scipy \
tensorboard \
transformers \
hf_transfer
transformers
CMD ["/bin/bash"]

View File

@@ -45,7 +45,6 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
numpy==1.26.4 \
scipy \
tensorboard \
transformers \
hf_transfer
transformers
CMD ["/bin/bash"]

View File

@@ -43,7 +43,6 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
numpy==1.26.4 \
scipy \
tensorboard \
transformers \
hf_transfer
transformers
CMD ["/bin/bash"]

View File

@@ -38,13 +38,11 @@ RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
datasets \
hf-doc-builder \
huggingface-hub \
hf_transfer \
Jinja2 \
librosa \
numpy==1.26.4 \
scipy \
tensorboard \
transformers \
hf_transfer
transformers
CMD ["/bin/bash"]

View File

@@ -16,35 +16,33 @@ RUN apt install -y bash \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.10 \
python3.10-dev \
python3.9 \
python3.9-dev \
python3-pip \
python3.10-venv && \
python3.9-venv && \
rm -rf /var/lib/apt/lists
# make sure to use venv
RUN python3.10 -m venv /opt/venv
RUN python3.9 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3.10 -m uv pip install --no-cache-dir \
RUN python3.9 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3.9 -m uv pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark && \
python3.10 -m pip install --no-cache-dir \
python3.9 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
huggingface-hub \
hf_transfer \
Jinja2 \
librosa \
numpy==1.26.4 \
scipy \
tensorboard \
transformers \
hf_transfer
transformers
CMD ["/bin/bash"]

View File

@@ -16,7 +16,6 @@ RUN apt install -y bash \
ca-certificates \
libsndfile1-dev \
python3.10 \
python3.10-dev \
python3-pip \
libgl1 \
python3.10-venv && \
@@ -44,7 +43,6 @@ RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
numpy==1.26.4 \
scipy \
tensorboard \
transformers matplotlib \
hf_transfer
transformers matplotlib
CMD ["/bin/bash"]

View File

@@ -17,7 +17,6 @@ RUN apt install -y bash \
libsndfile1-dev \
libgl1 \
python3.10 \
python3.10-dev \
python3-pip \
python3.10-venv && \
rm -rf /var/lib/apt/lists
@@ -38,14 +37,12 @@ RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
datasets \
hf-doc-builder \
huggingface-hub \
hf_transfer \
Jinja2 \
librosa \
numpy==1.26.4 \
scipy \
tensorboard \
transformers \
pytorch-lightning \
hf_transfer
pytorch-lightning
CMD ["/bin/bash"]

View File

@@ -17,7 +17,6 @@ RUN apt install -y bash \
libsndfile1-dev \
libgl1 \
python3.10 \
python3.10-dev \
python3-pip \
python3.10-venv && \
rm -rf /var/lib/apt/lists
@@ -38,14 +37,12 @@ RUN python3.10 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
datasets \
hf-doc-builder \
huggingface-hub \
hf_transfer \
Jinja2 \
librosa \
numpy==1.26.4 \
scipy \
tensorboard \
transformers \
xformers \
hf_transfer
xformers
CMD ["/bin/bash"]

View File

@@ -55,10 +55,8 @@
- sections:
- local: using-diffusers/overview_techniques
title: Overview
- local: using-diffusers/create_a_server
title: Create a server
- local: training/distributed_inference
title: Distributed inference
title: Distributed inference with multiple GPUs
- local: using-diffusers/merge_loras
title: Merge LoRAs
- local: using-diffusers/scheduler_features
@@ -77,8 +75,6 @@
title: Outpainting
title: Advanced inference
- sections:
- local: using-diffusers/cogvideox
title: CogVideoX
- local: using-diffusers/sdxl
title: Stable Diffusion XL
- local: using-diffusers/sdxl_turbo
@@ -133,8 +129,6 @@
title: T2I-Adapters
- local: training/instructpix2pix
title: InstructPix2Pix
- local: training/cogvideox
title: CogVideoX
title: Models
- isExpanded: false
sections:
@@ -152,16 +146,6 @@
title: Reinforcement learning training with DDPO
title: Methods
title: Training
- sections:
- local: quantization/overview
title: Getting Started
- local: quantization/bitsandbytes
title: bitsandbytes
- local: quantization/gguf
title: gguf
- local: quantization/torchao
title: torchao
title: Quantization Methods
- sections:
- local: optimization/fp16
title: Speed up inference
@@ -177,8 +161,6 @@
title: DeepCache
- local: optimization/tgate
title: TGATE
- local: optimization/xdit
title: xDiT
- sections:
- local: using-diffusers/stable_diffusion_jax_how_to
title: JAX/Flax
@@ -194,8 +176,6 @@
title: Metal Performance Shaders (MPS)
- local: optimization/habana
title: Habana Gaudi
- local: optimization/neuron
title: AWS Neuron
title: Optimized hardware
title: Accelerate inference and reduce memory
- sections:
@@ -210,10 +190,6 @@
- local: conceptual/evaluation
title: Evaluating Diffusion Models
title: Conceptual Guides
- sections:
- local: community_projects
title: Projects built with Diffusers
title: Community Projects
- sections:
- isExpanded: false
sections:
@@ -223,8 +199,6 @@
title: Logging
- local: api/outputs
title: Outputs
- local: api/quantization
title: Quantization
title: Main Classes
- isExpanded: false
sections:
@@ -245,109 +219,59 @@
sections:
- local: api/models/overview
title: Overview
- sections:
- local: api/models/controlnet
title: ControlNetModel
- local: api/models/controlnet_flux
title: FluxControlNetModel
- local: api/models/controlnet_hunyuandit
title: HunyuanDiT2DControlNetModel
- local: api/models/controlnet_sd3
title: SD3ControlNetModel
- local: api/models/controlnet_sparsectrl
title: SparseControlNetModel
- local: api/models/controlnet_union
title: ControlNetUnionModel
title: ControlNets
- sections:
- local: api/models/allegro_transformer3d
title: AllegroTransformer3DModel
- local: api/models/aura_flow_transformer2d
title: AuraFlowTransformer2DModel
- local: api/models/cogvideox_transformer3d
title: CogVideoXTransformer3DModel
- local: api/models/cogview3plus_transformer2d
title: CogView3PlusTransformer2DModel
- local: api/models/dit_transformer2d
title: DiTTransformer2DModel
- local: api/models/flux_transformer
title: FluxTransformer2DModel
- local: api/models/hunyuan_transformer2d
title: HunyuanDiT2DModel
- local: api/models/hunyuan_video_transformer_3d
title: HunyuanVideoTransformer3DModel
- local: api/models/latte_transformer3d
title: LatteTransformer3DModel
- local: api/models/lumina_nextdit2d
title: LuminaNextDiT2DModel
- local: api/models/ltx_video_transformer3d
title: LTXVideoTransformer3DModel
- local: api/models/mochi_transformer3d
title: MochiTransformer3DModel
- local: api/models/pixart_transformer2d
title: PixArtTransformer2DModel
- local: api/models/prior_transformer
title: PriorTransformer
- local: api/models/sd3_transformer2d
title: SD3Transformer2DModel
- local: api/models/sana_transformer2d
title: SanaTransformer2DModel
- local: api/models/stable_audio_transformer
title: StableAudioDiTModel
- local: api/models/transformer2d
title: Transformer2DModel
- local: api/models/transformer_temporal
title: TransformerTemporalModel
title: Transformers
- sections:
- local: api/models/stable_cascade_unet
title: StableCascadeUNet
- local: api/models/unet
title: UNet1DModel
- local: api/models/unet2d
title: UNet2DModel
- local: api/models/unet2d-cond
title: UNet2DConditionModel
- local: api/models/unet3d-cond
title: UNet3DConditionModel
- local: api/models/unet-motion
title: UNetMotionModel
- local: api/models/uvit2d
title: UViT2DModel
title: UNets
- sections:
- local: api/models/autoencoderkl
title: AutoencoderKL
- local: api/models/autoencoderkl_allegro
title: AutoencoderKLAllegro
- local: api/models/autoencoderkl_cogvideox
title: AutoencoderKLCogVideoX
- local: api/models/autoencoder_kl_hunyuan_video
title: AutoencoderKLHunyuanVideo
- local: api/models/autoencoderkl_ltx_video
title: AutoencoderKLLTXVideo
- local: api/models/autoencoderkl_mochi
title: AutoencoderKLMochi
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_dc
title: AutoencoderDC
- local: api/models/consistency_decoder_vae
title: ConsistencyDecoderVAE
- local: api/models/autoencoder_oobleck
title: Oobleck AutoEncoder
- local: api/models/autoencoder_tiny
title: Tiny AutoEncoder
- local: api/models/vq
title: VQModel
title: VAEs
- local: api/models/unet
title: UNet1DModel
- local: api/models/unet2d
title: UNet2DModel
- local: api/models/unet2d-cond
title: UNet2DConditionModel
- local: api/models/unet3d-cond
title: UNet3DConditionModel
- local: api/models/unet-motion
title: UNetMotionModel
- local: api/models/uvit2d
title: UViT2DModel
- local: api/models/vq
title: VQModel
- local: api/models/autoencoderkl
title: AutoencoderKL
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_tiny
title: Tiny AutoEncoder
- local: api/models/consistency_decoder_vae
title: ConsistencyDecoderVAE
- local: api/models/transformer2d
title: Transformer2DModel
- local: api/models/pixart_transformer2d
title: PixArtTransformer2DModel
- local: api/models/dit_transformer2d
title: DiTTransformer2DModel
- local: api/models/hunyuan_transformer2d
title: HunyuanDiT2DModel
- local: api/models/aura_flow_transformer2d
title: AuraFlowTransformer2DModel
- local: api/models/latte_transformer3d
title: LatteTransformer3DModel
- local: api/models/lumina_nextdit2d
title: LuminaNextDiT2DModel
- local: api/models/transformer_temporal
title: TransformerTemporalModel
- local: api/models/sd3_transformer2d
title: SD3Transformer2DModel
- local: api/models/prior_transformer
title: PriorTransformer
- local: api/models/controlnet
title: ControlNetModel
- local: api/models/controlnet_hunyuandit
title: HunyuanDiT2DControlNetModel
- local: api/models/controlnet_sd3
title: SD3ControlNetModel
title: Models
- isExpanded: false
sections:
- local: api/pipelines/overview
title: Overview
- local: api/pipelines/allegro
title: Allegro
- local: api/pipelines/amused
title: aMUSEd
- local: api/pipelines/animatediff
@@ -364,16 +288,10 @@
title: AutoPipeline
- local: api/pipelines/blip_diffusion
title: BLIP-Diffusion
- local: api/pipelines/cogvideox
title: CogVideoX
- local: api/pipelines/cogview3
title: CogView3
- local: api/pipelines/consistency_models
title: Consistency Models
- local: api/pipelines/controlnet
title: ControlNet
- local: api/pipelines/controlnet_flux
title: ControlNet with Flux.1
- local: api/pipelines/controlnet_hunyuandit
title: ControlNet with Hunyuan-DiT
- local: api/pipelines/controlnet_sd3
@@ -384,8 +302,6 @@
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
title: ControlNet-XS with Stable Diffusion XL
- local: api/pipelines/controlnet_union
title: ControlNetUnion
- local: api/pipelines/dance_diffusion
title: Dance Diffusion
- local: api/pipelines/ddim
@@ -398,12 +314,8 @@
title: DiffEdit
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/flux
title: Flux
- local: api/pipelines/hunyuandit
title: Hunyuan-DiT
- local: api/pipelines/hunyuan_video
title: HunyuanVideo
- local: api/pipelines/i2vgenxl
title: I2VGen-XL
- local: api/pipelines/pix2pix
@@ -420,18 +332,12 @@
title: Latent Consistency Models
- local: api/pipelines/latent_diffusion
title: Latent Diffusion
- local: api/pipelines/latte
title: Latte
- local: api/pipelines/ledits_pp
title: LEDITS++
- local: api/pipelines/ltx_video
title: LTX
- local: api/pipelines/lumina
title: Lumina-T2X
- local: api/pipelines/marigold
title: Marigold
- local: api/pipelines/mochi
title: Mochi
- local: api/pipelines/panorama
title: MultiDiffusion
- local: api/pipelines/musicldm
@@ -446,16 +352,12 @@
title: PixArt-α
- local: api/pipelines/pixart_sigma
title: PixArt-Σ
- local: api/pipelines/sana
title: Sana
- local: api/pipelines/self_attention_guidance
title: Self-Attention Guidance
- local: api/pipelines/semantic_stable_diffusion
title: Semantic Guidance
- local: api/pipelines/shap_e
title: Shap-E
- local: api/pipelines/stable_audio
title: Stable Audio
- local: api/pipelines/stable_cascade
title: Stable Cascade
- sections:
@@ -519,8 +421,6 @@
title: CMStochasticIterativeScheduler
- local: api/schedulers/consistency_decoder
title: ConsistencyDecoderScheduler
- local: api/schedulers/cosine_dpm
title: CosineDPMSolverMultistepScheduler
- local: api/schedulers/ddim_inverse
title: DDIMInverseScheduler
- local: api/schedulers/ddim

View File

@@ -15,133 +15,40 @@ specific language governing permissions and limitations under the License.
An attention processor is a class for applying different types of attention mechanisms.
## AttnProcessor
[[autodoc]] models.attention_processor.AttnProcessor
## AttnProcessor2_0
[[autodoc]] models.attention_processor.AttnProcessor2_0
## AttnAddedKVProcessor
[[autodoc]] models.attention_processor.AttnAddedKVProcessor
## AttnAddedKVProcessor2_0
[[autodoc]] models.attention_processor.AttnAddedKVProcessor2_0
[[autodoc]] models.attention_processor.AttnProcessorNPU
[[autodoc]] models.attention_processor.FusedAttnProcessor2_0
## Allegro
[[autodoc]] models.attention_processor.AllegroAttnProcessor2_0
## AuraFlow
[[autodoc]] models.attention_processor.AuraFlowAttnProcessor2_0
[[autodoc]] models.attention_processor.FusedAuraFlowAttnProcessor2_0
## CogVideoX
[[autodoc]] models.attention_processor.CogVideoXAttnProcessor2_0
[[autodoc]] models.attention_processor.FusedCogVideoXAttnProcessor2_0
## CrossFrameAttnProcessor
[[autodoc]] pipelines.text_to_video_synthesis.pipeline_text_to_video_zero.CrossFrameAttnProcessor
## Custom Diffusion
## CustomDiffusionAttnProcessor
[[autodoc]] models.attention_processor.CustomDiffusionAttnProcessor
## CustomDiffusionAttnProcessor2_0
[[autodoc]] models.attention_processor.CustomDiffusionAttnProcessor2_0
## CustomDiffusionXFormersAttnProcessor
[[autodoc]] models.attention_processor.CustomDiffusionXFormersAttnProcessor
## Flux
[[autodoc]] models.attention_processor.FluxAttnProcessor2_0
[[autodoc]] models.attention_processor.FusedFluxAttnProcessor2_0
[[autodoc]] models.attention_processor.FluxSingleAttnProcessor2_0
## Hunyuan
[[autodoc]] models.attention_processor.HunyuanAttnProcessor2_0
[[autodoc]] models.attention_processor.FusedHunyuanAttnProcessor2_0
[[autodoc]] models.attention_processor.PAGHunyuanAttnProcessor2_0
[[autodoc]] models.attention_processor.PAGCFGHunyuanAttnProcessor2_0
## IdentitySelfAttnProcessor2_0
[[autodoc]] models.attention_processor.PAGIdentitySelfAttnProcessor2_0
[[autodoc]] models.attention_processor.PAGCFGIdentitySelfAttnProcessor2_0
## IP-Adapter
[[autodoc]] models.attention_processor.IPAdapterAttnProcessor
[[autodoc]] models.attention_processor.IPAdapterAttnProcessor2_0
## JointAttnProcessor2_0
[[autodoc]] models.attention_processor.JointAttnProcessor2_0
[[autodoc]] models.attention_processor.PAGJointAttnProcessor2_0
[[autodoc]] models.attention_processor.PAGCFGJointAttnProcessor2_0
[[autodoc]] models.attention_processor.FusedJointAttnProcessor2_0
## LoRA
[[autodoc]] models.attention_processor.LoRAAttnProcessor
[[autodoc]] models.attention_processor.LoRAAttnProcessor2_0
[[autodoc]] models.attention_processor.LoRAAttnAddedKVProcessor
[[autodoc]] models.attention_processor.LoRAXFormersAttnProcessor
## Lumina-T2X
[[autodoc]] models.attention_processor.LuminaAttnProcessor2_0
## Mochi
[[autodoc]] models.attention_processor.MochiAttnProcessor2_0
[[autodoc]] models.attention_processor.MochiVaeAttnProcessor2_0
## Sana
[[autodoc]] models.attention_processor.SanaLinearAttnProcessor2_0
[[autodoc]] models.attention_processor.SanaMultiscaleAttnProcessor2_0
[[autodoc]] models.attention_processor.PAGCFGSanaLinearAttnProcessor2_0
[[autodoc]] models.attention_processor.PAGIdentitySanaLinearAttnProcessor2_0
## Stable Audio
[[autodoc]] models.attention_processor.StableAudioAttnProcessor2_0
## FusedAttnProcessor2_0
[[autodoc]] models.attention_processor.FusedAttnProcessor2_0
## SlicedAttnProcessor
[[autodoc]] models.attention_processor.SlicedAttnProcessor
## SlicedAttnAddedKVProcessor
[[autodoc]] models.attention_processor.SlicedAttnAddedKVProcessor
## XFormersAttnProcessor
[[autodoc]] models.attention_processor.XFormersAttnProcessor
[[autodoc]] models.attention_processor.XFormersAttnAddedKVProcessor
## XLAFlashAttnProcessor2_0
[[autodoc]] models.attention_processor.XLAFlashAttnProcessor2_0
## AttnProcessorNPU
[[autodoc]] models.attention_processor.AttnProcessorNPU

View File

@@ -12,16 +12,10 @@ specific language governing permissions and limitations under the License.
# LoRA
LoRA is a fast and lightweight training method that inserts and trains a significantly smaller number of parameters instead of all the model parameters. This produces a smaller file (~100 MBs) and makes it easier to quickly train a model to learn a new concept. LoRA weights are typically loaded into the denoiser, text encoder or both. The denoiser usually corresponds to a UNet ([`UNet2DConditionModel`], for example) or a Transformer ([`SD3Transformer2DModel`], for example). There are several classes for loading LoRA weights:
LoRA is a fast and lightweight training method that inserts and trains a significantly smaller number of parameters instead of all the model parameters. This produces a smaller file (~100 MBs) and makes it easier to quickly train a model to learn a new concept. LoRA weights are typically loaded into the UNet, text encoder or both. There are two classes for loading LoRA weights:
- [`StableDiffusionLoraLoaderMixin`] provides functions for loading and unloading, fusing and unfusing, enabling and disabling, and more functions for managing LoRA weights. This class can be used with any model.
- [`StableDiffusionXLLoraLoaderMixin`] is a [Stable Diffusion (SDXL)](../../api/pipelines/stable_diffusion/stable_diffusion_xl) version of the [`StableDiffusionLoraLoaderMixin`] class for loading and saving LoRA weights. It can only be used with the SDXL model.
- [`SD3LoraLoaderMixin`] provides similar functions for [Stable Diffusion 3](https://huggingface.co/blog/sd3).
- [`FluxLoraLoaderMixin`] provides similar functions for [Flux](https://huggingface.co/docs/diffusers/main/en/api/pipelines/flux).
- [`CogVideoXLoraLoaderMixin`] provides similar functions for [CogVideoX](https://huggingface.co/docs/diffusers/main/en/api/pipelines/cogvideox).
- [`Mochi1LoraLoaderMixin`] provides similar functions for [Mochi](https://huggingface.co/docs/diffusers/main/en/api/pipelines/mochi).
- [`AmusedLoraLoaderMixin`] is for the [`AmusedPipeline`].
- [`LoraBaseMixin`] provides a base class with several utility methods to fuse, unfuse, unload, LoRAs and more.
- [`LoraLoaderMixin`] provides functions for loading and unloading, fusing and unfusing, enabling and disabling, and more functions for managing LoRA weights. This class can be used with any model.
- [`StableDiffusionXLLoraLoaderMixin`] is a [Stable Diffusion (SDXL)](../../api/pipelines/stable_diffusion/stable_diffusion_xl) version of the [`LoraLoaderMixin`] class for loading and saving LoRA weights. It can only be used with the SDXL model.
<Tip>
@@ -29,34 +23,10 @@ To learn more about how to load LoRA weights, see the [LoRA](../../using-diffuse
</Tip>
## StableDiffusionLoraLoaderMixin
## LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.StableDiffusionLoraLoaderMixin
[[autodoc]] loaders.lora.LoraLoaderMixin
## StableDiffusionXLLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.StableDiffusionXLLoraLoaderMixin
## SD3LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.SD3LoraLoaderMixin
## FluxLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.FluxLoraLoaderMixin
## CogVideoXLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.CogVideoXLoraLoaderMixin
## Mochi1LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.Mochi1LoraLoaderMixin
## AmusedLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.AmusedLoraLoaderMixin
## LoraBaseMixin
[[autodoc]] loaders.lora_base.LoraBaseMixin
[[autodoc]] loaders.lora.StableDiffusionXLLoraLoaderMixin

View File

@@ -12,7 +12,7 @@ specific language governing permissions and limitations under the License.
# PEFT
Diffusers supports loading adapters such as [LoRA](../../using-diffusers/loading_adapters) with the [PEFT](https://huggingface.co/docs/peft/index) library with the [`~loaders.peft.PeftAdapterMixin`] class. This allows modeling classes in Diffusers like [`UNet2DConditionModel`], [`SD3Transformer2DModel`] to operate with an adapter.
Diffusers supports loading adapters such as [LoRA](../../using-diffusers/loading_adapters) with the [PEFT](https://huggingface.co/docs/peft/index) library with the [`~loaders.peft.PeftAdapterMixin`] class. This allows modeling classes in Diffusers like [`UNet2DConditionModel`] to load an adapter.
<Tip>

View File

@@ -51,7 +51,6 @@ The [`~loaders.FromSingleFileMixin.from_single_file`] method allows you to load:
- [`AutoencoderKL`]
- [`ControlNetModel`]
- [`SD3Transformer2DModel`]
- [`FluxTransformer2DModel`]
## FromSingleFileMixin

View File

@@ -12,7 +12,7 @@ specific language governing permissions and limitations under the License.
# UNet
Some training methods - like LoRA and Custom Diffusion - typically target the UNet's attention layers, but these training methods can also target other non-attention layers. Instead of training all of a model's parameters, only a subset of the parameters are trained, which is faster and more efficient. This class is useful if you're *only* loading weights into a UNet. If you need to load weights into the text encoder or a text encoder and UNet, try using the [`~loaders.StableDiffusionLoraLoaderMixin.load_lora_weights`] function instead.
Some training methods - like LoRA and Custom Diffusion - typically target the UNet's attention layers, but these training methods can also target other non-attention layers. Instead of training all of a model's parameters, only a subset of the parameters are trained, which is faster and more efficient. This class is useful if you're *only* loading weights into a UNet. If you need to load weights into the text encoder or a text encoder and UNet, try using the [`~loaders.LoraLoaderMixin.load_lora_weights`] function instead.
The [`UNet2DConditionLoadersMixin`] class provides functions for loading and saving weights, fusing and unfusing LoRAs, disabling and enabling LoRAs, and setting and deleting adapters.

View File

@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AllegroTransformer3DModel
A Diffusion Transformer model for 3D data from [Allegro](https://github.com/rhymes-ai/Allegro) was introduced in [Allegro: Open the Black Box of Commercial-Level Video Generation Model](https://huggingface.co/papers/2410.15458) by RhymesAI.
The model can be loaded with the following code snippet.
```python
from diffusers import AllegroTransformer3DModel
vae = AllegroTransformer3DModel.from_pretrained("rhymes-ai/Allegro", subfolder="transformer", torch_dtype=torch.bfloat16).to("cuda")
```
## AllegroTransformer3DModel
[[autodoc]] AllegroTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,72 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderDC
The 2D Autoencoder model used in [SANA](https://huggingface.co/papers/2410.10629) and introduced in [DCAE](https://huggingface.co/papers/2410.10733) by authors Junyu Chen\*, Han Cai\*, Junsong Chen, Enze Xie, Shang Yang, Haotian Tang, Muyang Li, Yao Lu, Song Han from MIT HAN Lab.
The abstract from the paper is:
*We present Deep Compression Autoencoder (DC-AE), a new family of autoencoder models for accelerating high-resolution diffusion models. Existing autoencoder models have demonstrated impressive results at a moderate spatial compression ratio (e.g., 8x), but fail to maintain satisfactory reconstruction accuracy for high spatial compression ratios (e.g., 64x). We address this challenge by introducing two key techniques: (1) Residual Autoencoding, where we design our models to learn residuals based on the space-to-channel transformed features to alleviate the optimization difficulty of high spatial-compression autoencoders; (2) Decoupled High-Resolution Adaptation, an efficient decoupled three-phases training strategy for mitigating the generalization penalty of high spatial-compression autoencoders. With these designs, we improve the autoencoder's spatial compression ratio up to 128 while maintaining the reconstruction quality. Applying our DC-AE to latent diffusion models, we achieve significant speedup without accuracy drop. For example, on ImageNet 512x512, our DC-AE provides 19.1x inference speedup and 17.9x training speedup on H100 GPU for UViT-H while achieving a better FID, compared with the widely used SD-VAE-f8 autoencoder. Our code is available at [this https URL](https://github.com/mit-han-lab/efficientvit).*
The following DCAE models are released and supported in Diffusers.
| Diffusers format | Original format |
|:----------------:|:---------------:|
| [`mit-han-lab/dc-ae-f32c32-sana-1.0-diffusers`](https://huggingface.co/mit-han-lab/dc-ae-f32c32-sana-1.0-diffusers) | [`mit-han-lab/dc-ae-f32c32-sana-1.0`](https://huggingface.co/mit-han-lab/dc-ae-f32c32-sana-1.0)
| [`mit-han-lab/dc-ae-f32c32-in-1.0-diffusers`](https://huggingface.co/mit-han-lab/dc-ae-f32c32-in-1.0-diffusers) | [`mit-han-lab/dc-ae-f32c32-in-1.0`](https://huggingface.co/mit-han-lab/dc-ae-f32c32-in-1.0)
| [`mit-han-lab/dc-ae-f32c32-mix-1.0-diffusers`](https://huggingface.co/mit-han-lab/dc-ae-f32c32-mix-1.0-diffusers) | [`mit-han-lab/dc-ae-f32c32-mix-1.0`](https://huggingface.co/mit-han-lab/dc-ae-f32c32-mix-1.0)
| [`mit-han-lab/dc-ae-f64c128-in-1.0-diffusers`](https://huggingface.co/mit-han-lab/dc-ae-f64c128-in-1.0-diffusers) | [`mit-han-lab/dc-ae-f64c128-in-1.0`](https://huggingface.co/mit-han-lab/dc-ae-f64c128-in-1.0)
| [`mit-han-lab/dc-ae-f64c128-mix-1.0-diffusers`](https://huggingface.co/mit-han-lab/dc-ae-f64c128-mix-1.0-diffusers) | [`mit-han-lab/dc-ae-f64c128-mix-1.0`](https://huggingface.co/mit-han-lab/dc-ae-f64c128-mix-1.0)
| [`mit-han-lab/dc-ae-f128c512-in-1.0-diffusers`](https://huggingface.co/mit-han-lab/dc-ae-f128c512-in-1.0-diffusers) | [`mit-han-lab/dc-ae-f128c512-in-1.0`](https://huggingface.co/mit-han-lab/dc-ae-f128c512-in-1.0)
| [`mit-han-lab/dc-ae-f128c512-mix-1.0-diffusers`](https://huggingface.co/mit-han-lab/dc-ae-f128c512-mix-1.0-diffusers) | [`mit-han-lab/dc-ae-f128c512-mix-1.0`](https://huggingface.co/mit-han-lab/dc-ae-f128c512-mix-1.0)
This model was contributed by [lawrence-cj](https://github.com/lawrence-cj).
Load a model in Diffusers format with [`~ModelMixin.from_pretrained`].
```python
from diffusers import AutoencoderDC
ae = AutoencoderDC.from_pretrained("mit-han-lab/dc-ae-f32c32-sana-1.0-diffusers", torch_dtype=torch.float32).to("cuda")
```
## Load a model in Diffusers via `from_single_file`
```python
from difusers import AutoencoderDC
ckpt_path = "https://huggingface.co/mit-han-lab/dc-ae-f32c32-sana-1.0/blob/main/model.safetensors"
model = AutoencoderDC.from_single_file(ckpt_path)
```
The `AutoencoderDC` model has `in` and `mix` single file checkpoint variants that have matching checkpoint keys, but use different scaling factors. It is not possible for Diffusers to automatically infer the correct config file to use with the model based on just the checkpoint and will default to configuring the model using the `mix` variant config file. To override the automatically determined config, please use the `config` argument when using single file loading with `in` variant checkpoints.
```python
from diffusers import AutoencoderDC
ckpt_path = "https://huggingface.co/mit-han-lab/dc-ae-f128c512-in-1.0/blob/main/model.safetensors"
model = AutoencoderDC.from_single_file(ckpt_path, config="mit-han-lab/dc-ae-f128c512-in-1.0-diffusers")
```
## AutoencoderDC
[[autodoc]] AutoencoderDC
- encode
- decode
- all
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

View File

@@ -1,32 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLHunyuanVideo
The 3D variational autoencoder (VAE) model with KL loss used in [HunyuanVideo](https://github.com/Tencent/HunyuanVideo/), which was introduced in [HunyuanVideo: A Systematic Framework For Large Video Generative Models](https://huggingface.co/papers/2412.03603) by Tencent.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLHunyuanVideo
vae = AutoencoderKLHunyuanVideo.from_pretrained("tencent/HunyuanVideo", torch_dtype=torch.float16)
```
## AutoencoderKLHunyuanVideo
[[autodoc]] AutoencoderKLHunyuanVideo
- decode
- all
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

View File

@@ -1,38 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# AutoencoderOobleck
The Oobleck variational autoencoder (VAE) model with KL loss was introduced in [Stability-AI/stable-audio-tools](https://github.com/Stability-AI/stable-audio-tools) and [Stable Audio Open](https://huggingface.co/papers/2407.14358) by Stability AI. The model is used in 🤗 Diffusers to encode audio waveforms into latents and to decode latent representations into audio waveforms.
The abstract from the paper is:
*Open generative models are vitally important for the community, allowing for fine-tunes and serving as baselines when presenting new models. However, most current text-to-audio models are private and not accessible for artists and researchers to build upon. Here we describe the architecture and training process of a new open-weights text-to-audio model trained with Creative Commons data. Our evaluation shows that the model's performance is competitive with the state-of-the-art across various metrics. Notably, the reported FDopenl3 results (measuring the realism of the generations) showcase its potential for high-quality stereo sound synthesis at 44.1kHz.*
## AutoencoderOobleck
[[autodoc]] AutoencoderOobleck
- decode
- encode
- all
## OobleckDecoderOutput
[[autodoc]] models.autoencoders.autoencoder_oobleck.OobleckDecoderOutput
## OobleckDecoderOutput
[[autodoc]] models.autoencoders.autoencoder_oobleck.OobleckDecoderOutput
## AutoencoderOobleckOutput
[[autodoc]] models.autoencoders.autoencoder_oobleck.AutoencoderOobleckOutput

View File

@@ -1,37 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLAllegro
The 3D variational autoencoder (VAE) model with KL loss used in [Allegro](https://github.com/rhymes-ai/Allegro) was introduced in [Allegro: Open the Black Box of Commercial-Level Video Generation Model](https://huggingface.co/papers/2410.15458) by RhymesAI.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLAllegro
vae = AutoencoderKLCogVideoX.from_pretrained("rhymes-ai/Allegro", subfolder="vae", torch_dtype=torch.float32).to("cuda")
```
## AutoencoderKLAllegro
[[autodoc]] AutoencoderKLAllegro
- decode
- encode
- all
## AutoencoderKLOutput
[[autodoc]] models.autoencoders.autoencoder_kl.AutoencoderKLOutput
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

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@@ -1,37 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLCogVideoX
The 3D variational autoencoder (VAE) model with KL loss used in [CogVideoX](https://github.com/THUDM/CogVideo) was introduced in [CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer](https://github.com/THUDM/CogVideo/blob/main/resources/CogVideoX.pdf) by Tsinghua University & ZhipuAI.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLCogVideoX
vae = AutoencoderKLCogVideoX.from_pretrained("THUDM/CogVideoX-2b", subfolder="vae", torch_dtype=torch.float16).to("cuda")
```
## AutoencoderKLCogVideoX
[[autodoc]] AutoencoderKLCogVideoX
- decode
- encode
- all
## AutoencoderKLOutput
[[autodoc]] models.autoencoders.autoencoder_kl.AutoencoderKLOutput
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

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@@ -1,37 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLLTXVideo
The 3D variational autoencoder (VAE) model with KL loss used in [LTX](https://huggingface.co/Lightricks/LTX-Video) was introduced by Lightricks.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLLTXVideo
vae = AutoencoderKLLTXVideo.from_pretrained("TODO/TODO", subfolder="vae", torch_dtype=torch.float32).to("cuda")
```
## AutoencoderKLLTXVideo
[[autodoc]] AutoencoderKLLTXVideo
- decode
- encode
- all
## AutoencoderKLOutput
[[autodoc]] models.autoencoders.autoencoder_kl.AutoencoderKLOutput
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

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@@ -1,32 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLMochi
The 3D variational autoencoder (VAE) model with KL loss used in [Mochi](https://github.com/genmoai/models) was introduced in [Mochi 1 Preview](https://huggingface.co/genmo/mochi-1-preview) by Tsinghua University & ZhipuAI.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLMochi
vae = AutoencoderKLMochi.from_pretrained("genmo/mochi-1-preview", subfolder="vae", torch_dtype=torch.float32).to("cuda")
```
## AutoencoderKLMochi
[[autodoc]] AutoencoderKLMochi
- decode
- all
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

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@@ -1,30 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# CogVideoXTransformer3DModel
A Diffusion Transformer model for 3D data from [CogVideoX](https://github.com/THUDM/CogVideo) was introduced in [CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer](https://github.com/THUDM/CogVideo/blob/main/resources/CogVideoX.pdf) by Tsinghua University & ZhipuAI.
The model can be loaded with the following code snippet.
```python
from diffusers import CogVideoXTransformer3DModel
vae = CogVideoXTransformer3DModel.from_pretrained("THUDM/CogVideoX-2b", subfolder="transformer", torch_dtype=torch.float16).to("cuda")
```
## CogVideoXTransformer3DModel
[[autodoc]] CogVideoXTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

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@@ -1,30 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# CogView3PlusTransformer2DModel
A Diffusion Transformer model for 2D data from [CogView3Plus](https://github.com/THUDM/CogView3) was introduced in [CogView3: Finer and Faster Text-to-Image Generation via Relay Diffusion](https://huggingface.co/papers/2403.05121) by Tsinghua University & ZhipuAI.
The model can be loaded with the following code snippet.
```python
from diffusers import CogView3PlusTransformer2DModel
vae = CogView3PlusTransformer2DModel.from_pretrained("THUDM/CogView3Plus-3b", subfolder="transformer", torch_dtype=torch.bfloat16).to("cuda")
```
## CogView3PlusTransformer2DModel
[[autodoc]] CogView3PlusTransformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -29,7 +29,7 @@ from diffusers import StableDiffusionControlNetPipeline, ControlNetModel
url = "https://huggingface.co/lllyasviel/ControlNet-v1-1/blob/main/control_v11p_sd15_canny.pth" # can also be a local path
controlnet = ControlNetModel.from_single_file(url)
url = "https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5/blob/main/v1-5-pruned.safetensors" # can also be a local path
url = "https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/v1-5-pruned.safetensors" # can also be a local path
pipe = StableDiffusionControlNetPipeline.from_single_file(url, controlnet=controlnet)
```
@@ -39,7 +39,7 @@ pipe = StableDiffusionControlNetPipeline.from_single_file(url, controlnet=contro
## ControlNetOutput
[[autodoc]] models.controlnets.controlnet.ControlNetOutput
[[autodoc]] models.controlnet.ControlNetOutput
## FlaxControlNetModel
@@ -47,4 +47,4 @@ pipe = StableDiffusionControlNetPipeline.from_single_file(url, controlnet=contro
## FlaxControlNetOutput
[[autodoc]] models.controlnets.controlnet_flax.FlaxControlNetOutput
[[autodoc]] models.controlnet_flax.FlaxControlNetOutput

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@@ -1,45 +0,0 @@
<!--Copyright 2024 The HuggingFace Team and The InstantX Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# FluxControlNetModel
FluxControlNetModel is an implementation of ControlNet for Flux.1.
The ControlNet model was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, Maneesh Agrawala. It provides a greater degree of control over text-to-image generation by conditioning the model on additional inputs such as edge maps, depth maps, segmentation maps, and keypoints for pose detection.
The abstract from the paper is:
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
## Loading from the original format
By default the [`FluxControlNetModel`] should be loaded with [`~ModelMixin.from_pretrained`].
```py
from diffusers import FluxControlNetPipeline
from diffusers.models import FluxControlNetModel, FluxMultiControlNetModel
controlnet = FluxControlNetModel.from_pretrained("InstantX/FLUX.1-dev-Controlnet-Canny")
pipe = FluxControlNetPipeline.from_pretrained("black-forest-labs/FLUX.1-dev", controlnet=controlnet)
controlnet = FluxControlNetModel.from_pretrained("InstantX/FLUX.1-dev-Controlnet-Canny")
controlnet = FluxMultiControlNetModel([controlnet])
pipe = FluxControlNetPipeline.from_pretrained("black-forest-labs/FLUX.1-dev", controlnet=controlnet)
```
## FluxControlNetModel
[[autodoc]] FluxControlNetModel
## FluxControlNetOutput
[[autodoc]] models.controlnet_flux.FluxControlNetOutput

View File

@@ -38,5 +38,5 @@ pipe = StableDiffusion3ControlNetPipeline.from_pretrained("stabilityai/stable-di
## SD3ControlNetOutput
[[autodoc]] models.controlnets.controlnet_sd3.SD3ControlNetOutput
[[autodoc]] models.controlnet_sd3.SD3ControlNetOutput

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@@ -1,46 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# SparseControlNetModel
SparseControlNetModel is an implementation of ControlNet for [AnimateDiff](https://arxiv.org/abs/2307.04725).
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
The SparseCtrl version of ControlNet was introduced in [SparseCtrl: Adding Sparse Controls to Text-to-Video Diffusion Models](https://arxiv.org/abs/2311.16933) for achieving controlled generation in text-to-video diffusion models by Yuwei Guo, Ceyuan Yang, Anyi Rao, Maneesh Agrawala, Dahua Lin, and Bo Dai.
The abstract from the paper is:
*The development of text-to-video (T2V), i.e., generating videos with a given text prompt, has been significantly advanced in recent years. However, relying solely on text prompts often results in ambiguous frame composition due to spatial uncertainty. The research community thus leverages the dense structure signals, e.g., per-frame depth/edge sequences, to enhance controllability, whose collection accordingly increases the burden of inference. In this work, we present SparseCtrl to enable flexible structure control with temporally sparse signals, requiring only one or a few inputs, as shown in Figure 1. It incorporates an additional condition encoder to process these sparse signals while leaving the pre-trained T2V model untouched. The proposed approach is compatible with various modalities, including sketches, depth maps, and RGB images, providing more practical control for video generation and promoting applications such as storyboarding, depth rendering, keyframe animation, and interpolation. Extensive experiments demonstrate the generalization of SparseCtrl on both original and personalized T2V generators. Codes and models will be publicly available at [this https URL](https://guoyww.github.io/projects/SparseCtrl).*
## Example for loading SparseControlNetModel
```python
import torch
from diffusers import SparseControlNetModel
# fp32 variant in float16
# 1. Scribble checkpoint
controlnet = SparseControlNetModel.from_pretrained("guoyww/animatediff-sparsectrl-scribble", torch_dtype=torch.float16)
# 2. RGB checkpoint
controlnet = SparseControlNetModel.from_pretrained("guoyww/animatediff-sparsectrl-rgb", torch_dtype=torch.float16)
# For loading fp16 variant, pass `variant="fp16"` as an additional parameter
```
## SparseControlNetModel
[[autodoc]] SparseControlNetModel
## SparseControlNetOutput
[[autodoc]] models.controlnet_sparsectrl.SparseControlNetOutput

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@@ -1,35 +0,0 @@
<!--Copyright 2024 The HuggingFace Team and The InstantX Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNetUnionModel
ControlNetUnionModel is an implementation of ControlNet for Stable Diffusion XL.
The ControlNet model was introduced in [ControlNetPlus](https://github.com/xinsir6/ControlNetPlus) by xinsir6. It supports multiple conditioning inputs without increasing computation.
*We design a new architecture that can support 10+ control types in condition text-to-image generation and can generate high resolution images visually comparable with midjourney. The network is based on the original ControlNet architecture, we propose two new modules to: 1 Extend the original ControlNet to support different image conditions using the same network parameter. 2 Support multiple conditions input without increasing computation offload, which is especially important for designers who want to edit image in detail, different conditions use the same condition encoder, without adding extra computations or parameters.*
## Loading
By default the [`ControlNetUnionModel`] should be loaded with [`~ModelMixin.from_pretrained`].
```py
from diffusers import StableDiffusionXLControlNetUnionPipeline, ControlNetUnionModel
controlnet = ControlNetUnionModel.from_pretrained("xinsir/controlnet-union-sdxl-1.0")
pipe = StableDiffusionXLControlNetUnionPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", controlnet=controlnet)
```
## ControlNetUnionModel
[[autodoc]] ControlNetUnionModel

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@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# FluxTransformer2DModel
A Transformer model for image-like data from [Flux](https://blackforestlabs.ai/announcing-black-forest-labs/).
## FluxTransformer2DModel
[[autodoc]] FluxTransformer2DModel

View File

@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# HunyuanVideoTransformer3DModel
A Diffusion Transformer model for 3D video-like data was introduced in [HunyuanVideo: A Systematic Framework For Large Video Generative Models](https://huggingface.co/papers/2412.03603) by Tencent.
The model can be loaded with the following code snippet.
```python
from diffusers import HunyuanVideoTransformer3DModel
transformer = HunyuanVideoTransformer3DModel.from_pretrained("tencent/HunyuanVideo", torch_dtype=torch.bfloat16)
```
## HunyuanVideoTransformer3DModel
[[autodoc]] HunyuanVideoTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

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@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# LTXVideoTransformer3DModel
A Diffusion Transformer model for 3D data from [LTX](https://huggingface.co/Lightricks/LTX-Video) was introduced by Lightricks.
The model can be loaded with the following code snippet.
```python
from diffusers import LTXVideoTransformer3DModel
transformer = LTXVideoTransformer3DModel.from_pretrained("TODO/TODO", subfolder="transformer", torch_dtype=torch.bfloat16).to("cuda")
```
## LTXVideoTransformer3DModel
[[autodoc]] LTXVideoTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

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@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# MochiTransformer3DModel
A Diffusion Transformer model for 3D video-like data was introduced in [Mochi-1 Preview](https://huggingface.co/genmo/mochi-1-preview) by Genmo.
The model can be loaded with the following code snippet.
```python
from diffusers import MochiTransformer3DModel
vae = MochiTransformer3DModel.from_pretrained("genmo/mochi-1-preview", subfolder="transformer", torch_dtype=torch.float16).to("cuda")
```
## MochiTransformer3DModel
[[autodoc]] MochiTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

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@@ -1,34 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# SanaTransformer2DModel
A Diffusion Transformer model for 2D data from [SANA: Efficient High-Resolution Image Synthesis with Linear Diffusion Transformers](https://huggingface.co/papers/2410.10629) was introduced from NVIDIA and MIT HAN Lab, by Enze Xie, Junsong Chen, Junyu Chen, Han Cai, Haotian Tang, Yujun Lin, Zhekai Zhang, Muyang Li, Ligeng Zhu, Yao Lu, Song Han.
The abstract from the paper is:
*We introduce Sana, a text-to-image framework that can efficiently generate images up to 4096×4096 resolution. Sana can synthesize high-resolution, high-quality images with strong text-image alignment at a remarkably fast speed, deployable on laptop GPU. Core designs include: (1) Deep compression autoencoder: unlike traditional AEs, which compress images only 8×, we trained an AE that can compress images 32×, effectively reducing the number of latent tokens. (2) Linear DiT: we replace all vanilla attention in DiT with linear attention, which is more efficient at high resolutions without sacrificing quality. (3) Decoder-only text encoder: we replaced T5 with modern decoder-only small LLM as the text encoder and designed complex human instruction with in-context learning to enhance the image-text alignment. (4) Efficient training and sampling: we propose Flow-DPM-Solver to reduce sampling steps, with efficient caption labeling and selection to accelerate convergence. As a result, Sana-0.6B is very competitive with modern giant diffusion model (e.g. Flux-12B), being 20 times smaller and 100+ times faster in measured throughput. Moreover, Sana-0.6B can be deployed on a 16GB laptop GPU, taking less than 1 second to generate a 1024×1024 resolution image. Sana enables content creation at low cost. Code and model will be publicly released.*
The model can be loaded with the following code snippet.
```python
from diffusers import SanaTransformer2DModel
transformer = SanaTransformer2DModel.from_pretrained("Efficient-Large-Model/Sana_1600M_1024px_diffusers", subfolder="transformer", torch_dtype=torch.float16)
```
## SanaTransformer2DModel
[[autodoc]] SanaTransformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

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@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# StableAudioDiTModel
A Transformer model for audio waveforms from [Stable Audio Open](https://huggingface.co/papers/2407.14358).
## StableAudioDiTModel
[[autodoc]] StableAudioDiTModel

View File

@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# StableCascadeUNet
A UNet model from the [Stable Cascade pipeline](../pipelines/stable_cascade.md).
## StableCascadeUNet
[[autodoc]] models.unets.unet_stable_cascade.StableCascadeUNet

View File

@@ -1,34 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# Allegro
[Allegro: Open the Black Box of Commercial-Level Video Generation Model](https://huggingface.co/papers/2410.15458) from RhymesAI, by Yuan Zhou, Qiuyue Wang, Yuxuan Cai, Huan Yang.
The abstract from the paper is:
*Significant advancements have been made in the field of video generation, with the open-source community contributing a wealth of research papers and tools for training high-quality models. However, despite these efforts, the available information and resources remain insufficient for achieving commercial-level performance. In this report, we open the black box and introduce Allegro, an advanced video generation model that excels in both quality and temporal consistency. We also highlight the current limitations in the field and present a comprehensive methodology for training high-performance, commercial-level video generation models, addressing key aspects such as data, model architecture, training pipeline, and evaluation. Our user study shows that Allegro surpasses existing open-source models and most commercial models, ranking just behind Hailuo and Kling. Code: https://github.com/rhymes-ai/Allegro , Model: https://huggingface.co/rhymes-ai/Allegro , Gallery: https://rhymes.ai/allegro_gallery .*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## AllegroPipeline
[[autodoc]] AllegroPipeline
- all
- __call__
## AllegroPipelineOutput
[[autodoc]] pipelines.allegro.pipeline_output.AllegroPipelineOutput

View File

@@ -25,11 +25,7 @@ The abstract of the paper is the following:
| Pipeline | Tasks | Demo
|---|---|:---:|
| [AnimateDiffPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff.py) | *Text-to-Video Generation with AnimateDiff* |
| [AnimateDiffControlNetPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_controlnet.py) | *Controlled Video-to-Video Generation with AnimateDiff using ControlNet* |
| [AnimateDiffSparseControlNetPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_sparsectrl.py) | *Controlled Video-to-Video Generation with AnimateDiff using SparseCtrl* |
| [AnimateDiffSDXLPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_sdxl.py) | *Video-to-Video Generation with AnimateDiff* |
| [AnimateDiffVideoToVideoPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video.py) | *Video-to-Video Generation with AnimateDiff* |
| [AnimateDiffVideoToVideoControlNetPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/animatediff/pipeline_animatediff_video2video_controlnet.py) | *Video-to-Video Generation with AnimateDiff using ControlNet* |
## Available checkpoints
@@ -104,266 +100,6 @@ AnimateDiff tends to work better with finetuned Stable Diffusion models. If you
</Tip>
### AnimateDiffControlNetPipeline
AnimateDiff can also be used with ControlNets ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala. With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide depth maps, the ControlNet model generates a video that'll preserve the spatial information from the depth maps. It is a more flexible and accurate way to control the video generation process.
```python
import torch
from diffusers import AnimateDiffControlNetPipeline, AutoencoderKL, ControlNetModel, MotionAdapter, LCMScheduler
from diffusers.utils import export_to_gif, load_video
# Additionally, you will need a preprocess videos before they can be used with the ControlNet
# HF maintains just the right package for it: `pip install controlnet_aux`
from controlnet_aux.processor import ZoeDetector
# Download controlnets from https://huggingface.co/lllyasviel/ControlNet-v1-1 to use .from_single_file
# Download Diffusers-format controlnets, such as https://huggingface.co/lllyasviel/sd-controlnet-depth, to use .from_pretrained()
controlnet = ControlNetModel.from_single_file("control_v11f1p_sd15_depth.pth", torch_dtype=torch.float16)
# We use AnimateLCM for this example but one can use the original motion adapters as well (for example, https://huggingface.co/guoyww/animatediff-motion-adapter-v1-5-3)
motion_adapter = MotionAdapter.from_pretrained("wangfuyun/AnimateLCM")
vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse", torch_dtype=torch.float16)
pipe: AnimateDiffControlNetPipeline = AnimateDiffControlNetPipeline.from_pretrained(
"SG161222/Realistic_Vision_V5.1_noVAE",
motion_adapter=motion_adapter,
controlnet=controlnet,
vae=vae,
).to(device="cuda", dtype=torch.float16)
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config, beta_schedule="linear")
pipe.load_lora_weights("wangfuyun/AnimateLCM", weight_name="AnimateLCM_sd15_t2v_lora.safetensors", adapter_name="lcm-lora")
pipe.set_adapters(["lcm-lora"], [0.8])
depth_detector = ZoeDetector.from_pretrained("lllyasviel/Annotators").to("cuda")
video = load_video("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-vid2vid-input-1.gif")
conditioning_frames = []
with pipe.progress_bar(total=len(video)) as progress_bar:
for frame in video:
conditioning_frames.append(depth_detector(frame))
progress_bar.update()
prompt = "a panda, playing a guitar, sitting in a pink boat, in the ocean, mountains in background, realistic, high quality"
negative_prompt = "bad quality, worst quality"
video = pipe(
prompt=prompt,
negative_prompt=negative_prompt,
num_frames=len(video),
num_inference_steps=10,
guidance_scale=2.0,
conditioning_frames=conditioning_frames,
generator=torch.Generator().manual_seed(42),
).frames[0]
export_to_gif(video, "animatediff_controlnet.gif", fps=8)
```
Here are some sample outputs:
<table align="center">
<tr>
<th align="center">Source Video</th>
<th align="center">Output Video</th>
</tr>
<tr>
<td align="center">
raccoon playing a guitar
<br />
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-vid2vid-input-1.gif" alt="racoon playing a guitar" />
</td>
<td align="center">
a panda, playing a guitar, sitting in a pink boat, in the ocean, mountains in background, realistic, high quality
<br/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-controlnet-output.gif" alt="a panda, playing a guitar, sitting in a pink boat, in the ocean, mountains in background, realistic, high quality" />
</td>
</tr>
</table>
### AnimateDiffSparseControlNetPipeline
[SparseCtrl: Adding Sparse Controls to Text-to-Video Diffusion Models](https://arxiv.org/abs/2311.16933) for achieving controlled generation in text-to-video diffusion models by Yuwei Guo, Ceyuan Yang, Anyi Rao, Maneesh Agrawala, Dahua Lin, and Bo Dai.
The abstract from the paper is:
*The development of text-to-video (T2V), i.e., generating videos with a given text prompt, has been significantly advanced in recent years. However, relying solely on text prompts often results in ambiguous frame composition due to spatial uncertainty. The research community thus leverages the dense structure signals, e.g., per-frame depth/edge sequences, to enhance controllability, whose collection accordingly increases the burden of inference. In this work, we present SparseCtrl to enable flexible structure control with temporally sparse signals, requiring only one or a few inputs, as shown in Figure 1. It incorporates an additional condition encoder to process these sparse signals while leaving the pre-trained T2V model untouched. The proposed approach is compatible with various modalities, including sketches, depth maps, and RGB images, providing more practical control for video generation and promoting applications such as storyboarding, depth rendering, keyframe animation, and interpolation. Extensive experiments demonstrate the generalization of SparseCtrl on both original and personalized T2V generators. Codes and models will be publicly available at [this https URL](https://guoyww.github.io/projects/SparseCtrl).*
SparseCtrl introduces the following checkpoints for controlled text-to-video generation:
- [SparseCtrl Scribble](https://huggingface.co/guoyww/animatediff-sparsectrl-scribble)
- [SparseCtrl RGB](https://huggingface.co/guoyww/animatediff-sparsectrl-rgb)
#### Using SparseCtrl Scribble
```python
import torch
from diffusers import AnimateDiffSparseControlNetPipeline
from diffusers.models import AutoencoderKL, MotionAdapter, SparseControlNetModel
from diffusers.schedulers import DPMSolverMultistepScheduler
from diffusers.utils import export_to_gif, load_image
model_id = "SG161222/Realistic_Vision_V5.1_noVAE"
motion_adapter_id = "guoyww/animatediff-motion-adapter-v1-5-3"
controlnet_id = "guoyww/animatediff-sparsectrl-scribble"
lora_adapter_id = "guoyww/animatediff-motion-lora-v1-5-3"
vae_id = "stabilityai/sd-vae-ft-mse"
device = "cuda"
motion_adapter = MotionAdapter.from_pretrained(motion_adapter_id, torch_dtype=torch.float16).to(device)
controlnet = SparseControlNetModel.from_pretrained(controlnet_id, torch_dtype=torch.float16).to(device)
vae = AutoencoderKL.from_pretrained(vae_id, torch_dtype=torch.float16).to(device)
scheduler = DPMSolverMultistepScheduler.from_pretrained(
model_id,
subfolder="scheduler",
beta_schedule="linear",
algorithm_type="dpmsolver++",
use_karras_sigmas=True,
)
pipe = AnimateDiffSparseControlNetPipeline.from_pretrained(
model_id,
motion_adapter=motion_adapter,
controlnet=controlnet,
vae=vae,
scheduler=scheduler,
torch_dtype=torch.float16,
).to(device)
pipe.load_lora_weights(lora_adapter_id, adapter_name="motion_lora")
pipe.fuse_lora(lora_scale=1.0)
prompt = "an aerial view of a cyberpunk city, night time, neon lights, masterpiece, high quality"
negative_prompt = "low quality, worst quality, letterboxed"
image_files = [
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-1.png",
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-2.png",
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-3.png"
]
condition_frame_indices = [0, 8, 15]
conditioning_frames = [load_image(img_file) for img_file in image_files]
video = pipe(
prompt=prompt,
negative_prompt=negative_prompt,
num_inference_steps=25,
conditioning_frames=conditioning_frames,
controlnet_conditioning_scale=1.0,
controlnet_frame_indices=condition_frame_indices,
generator=torch.Generator().manual_seed(1337),
).frames[0]
export_to_gif(video, "output.gif")
```
Here are some sample outputs:
<table align="center">
<tr>
<center>
<b>an aerial view of a cyberpunk city, night time, neon lights, masterpiece, high quality</b>
</center>
</tr>
<tr>
<td>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-1.png" alt="scribble-1" />
</center>
</td>
<td>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-2.png" alt="scribble-2" />
</center>
</td>
<td>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-scribble-3.png" alt="scribble-3" />
</center>
</td>
</tr>
<tr>
<td colspan=3>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-sparsectrl-scribble-results.gif" alt="an aerial view of a cyberpunk city, night time, neon lights, masterpiece, high quality" />
</center>
</td>
</tr>
</table>
#### Using SparseCtrl RGB
```python
import torch
from diffusers import AnimateDiffSparseControlNetPipeline
from diffusers.models import AutoencoderKL, MotionAdapter, SparseControlNetModel
from diffusers.schedulers import DPMSolverMultistepScheduler
from diffusers.utils import export_to_gif, load_image
model_id = "SG161222/Realistic_Vision_V5.1_noVAE"
motion_adapter_id = "guoyww/animatediff-motion-adapter-v1-5-3"
controlnet_id = "guoyww/animatediff-sparsectrl-rgb"
lora_adapter_id = "guoyww/animatediff-motion-lora-v1-5-3"
vae_id = "stabilityai/sd-vae-ft-mse"
device = "cuda"
motion_adapter = MotionAdapter.from_pretrained(motion_adapter_id, torch_dtype=torch.float16).to(device)
controlnet = SparseControlNetModel.from_pretrained(controlnet_id, torch_dtype=torch.float16).to(device)
vae = AutoencoderKL.from_pretrained(vae_id, torch_dtype=torch.float16).to(device)
scheduler = DPMSolverMultistepScheduler.from_pretrained(
model_id,
subfolder="scheduler",
beta_schedule="linear",
algorithm_type="dpmsolver++",
use_karras_sigmas=True,
)
pipe = AnimateDiffSparseControlNetPipeline.from_pretrained(
model_id,
motion_adapter=motion_adapter,
controlnet=controlnet,
vae=vae,
scheduler=scheduler,
torch_dtype=torch.float16,
).to(device)
pipe.load_lora_weights(lora_adapter_id, adapter_name="motion_lora")
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-firework.png")
video = pipe(
prompt="closeup face photo of man in black clothes, night city street, bokeh, fireworks in background",
negative_prompt="low quality, worst quality",
num_inference_steps=25,
conditioning_frames=image,
controlnet_frame_indices=[0],
controlnet_conditioning_scale=1.0,
generator=torch.Generator().manual_seed(42),
).frames[0]
export_to_gif(video, "output.gif")
```
Here are some sample outputs:
<table align="center">
<tr>
<center>
<b>closeup face photo of man in black clothes, night city street, bokeh, fireworks in background</b>
</center>
</tr>
<tr>
<td>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-firework.png" alt="closeup face photo of man in black clothes, night city street, bokeh, fireworks in background" />
</center>
</td>
<td>
<center>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-sparsectrl-rgb-result.gif" alt="closeup face photo of man in black clothes, night city street, bokeh, fireworks in background" />
</center>
</td>
</tr>
</table>
### AnimateDiffSDXLPipeline
AnimateDiff can also be used with SDXL models. This is currently an experimental feature as only a beta release of the motion adapter checkpoint is available.
@@ -519,97 +255,6 @@ Here are some sample outputs:
</tr>
</table>
### AnimateDiffVideoToVideoControlNetPipeline
AnimateDiff can be used together with ControlNets to enhance video-to-video generation by allowing for precise control over the output. ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala, and allows you to condition Stable Diffusion with an additional control image to ensure that the spatial information is preserved throughout the video.
This pipeline allows you to condition your generation both on the original video and on a sequence of control images.
```python
import torch
from PIL import Image
from tqdm.auto import tqdm
from controlnet_aux.processor import OpenposeDetector
from diffusers import AnimateDiffVideoToVideoControlNetPipeline
from diffusers.utils import export_to_gif, load_video
from diffusers import AutoencoderKL, ControlNetModel, MotionAdapter, LCMScheduler
# Load the ControlNet
controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-openpose", torch_dtype=torch.float16)
# Load the motion adapter
motion_adapter = MotionAdapter.from_pretrained("wangfuyun/AnimateLCM")
# Load SD 1.5 based finetuned model
vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse", torch_dtype=torch.float16)
pipe = AnimateDiffVideoToVideoControlNetPipeline.from_pretrained(
"SG161222/Realistic_Vision_V5.1_noVAE",
motion_adapter=motion_adapter,
controlnet=controlnet,
vae=vae,
).to(device="cuda", dtype=torch.float16)
# Enable LCM to speed up inference
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config, beta_schedule="linear")
pipe.load_lora_weights("wangfuyun/AnimateLCM", weight_name="AnimateLCM_sd15_t2v_lora.safetensors", adapter_name="lcm-lora")
pipe.set_adapters(["lcm-lora"], [0.8])
video = load_video("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/dance.gif")
video = [frame.convert("RGB") for frame in video]
prompt = "astronaut in space, dancing"
negative_prompt = "bad quality, worst quality, jpeg artifacts, ugly"
# Create controlnet preprocessor
open_pose = OpenposeDetector.from_pretrained("lllyasviel/Annotators").to("cuda")
# Preprocess controlnet images
conditioning_frames = []
for frame in tqdm(video):
conditioning_frames.append(open_pose(frame))
strength = 0.8
with torch.inference_mode():
video = pipe(
video=video,
prompt=prompt,
negative_prompt=negative_prompt,
num_inference_steps=10,
guidance_scale=2.0,
controlnet_conditioning_scale=0.75,
conditioning_frames=conditioning_frames,
strength=strength,
generator=torch.Generator().manual_seed(42),
).frames[0]
video = [frame.resize(conditioning_frames[0].size) for frame in video]
export_to_gif(video, f"animatediff_vid2vid_controlnet.gif", fps=8)
```
Here are some sample outputs:
<table align="center">
<tr>
<th align="center">Source Video</th>
<th align="center">Output Video</th>
</tr>
<tr>
<td align="center">
anime girl, dancing
<br />
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/dance.gif" alt="anime girl, dancing" />
</td>
<td align="center">
astronaut in space, dancing
<br/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff_vid2vid_controlnet.gif" alt="astronaut in space, dancing" />
</td>
</tr>
</table>
**The lights and composition were transferred from the Source Video.**
## Using Motion LoRAs
Motion LoRAs are a collection of LoRAs that work with the `guoyww/animatediff-motion-adapter-v1-5-2` checkpoint. These LoRAs are responsible for adding specific types of motion to the animations.
@@ -914,89 +559,6 @@ export_to_gif(frames, "animatelcm-motion-lora.gif")
</tr>
</table>
## Using FreeNoise
[FreeNoise: Tuning-Free Longer Video Diffusion via Noise Rescheduling](https://arxiv.org/abs/2310.15169) by Haonan Qiu, Menghan Xia, Yong Zhang, Yingqing He, Xintao Wang, Ying Shan, Ziwei Liu.
FreeNoise is a sampling mechanism that can generate longer videos with short-video generation models by employing noise-rescheduling, temporal attention over sliding windows, and weighted averaging of latent frames. It also can be used with multiple prompts to allow for interpolated video generations. More details are available in the paper.
The currently supported AnimateDiff pipelines that can be used with FreeNoise are:
- [`AnimateDiffPipeline`]
- [`AnimateDiffControlNetPipeline`]
- [`AnimateDiffVideoToVideoPipeline`]
- [`AnimateDiffVideoToVideoControlNetPipeline`]
In order to use FreeNoise, a single line needs to be added to the inference code after loading your pipelines.
```diff
+ pipe.enable_free_noise()
```
After this, either a single prompt could be used, or multiple prompts can be passed as a dictionary of integer-string pairs. The integer keys of the dictionary correspond to the frame index at which the influence of that prompt would be maximum. Each frame index should map to a single string prompt. The prompts for intermediate frame indices, that are not passed in the dictionary, are created by interpolating between the frame prompts that are passed. By default, simple linear interpolation is used. However, you can customize this behaviour with a callback to the `prompt_interpolation_callback` parameter when enabling FreeNoise.
Full example:
```python
import torch
from diffusers import AutoencoderKL, AnimateDiffPipeline, LCMScheduler, MotionAdapter
from diffusers.utils import export_to_video, load_image
# Load pipeline
dtype = torch.float16
motion_adapter = MotionAdapter.from_pretrained("wangfuyun/AnimateLCM", torch_dtype=dtype)
vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse", torch_dtype=dtype)
pipe = AnimateDiffPipeline.from_pretrained("emilianJR/epiCRealism", motion_adapter=motion_adapter, vae=vae, torch_dtype=dtype)
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config, beta_schedule="linear")
pipe.load_lora_weights(
"wangfuyun/AnimateLCM", weight_name="AnimateLCM_sd15_t2v_lora.safetensors", adapter_name="lcm_lora"
)
pipe.set_adapters(["lcm_lora"], [0.8])
# Enable FreeNoise for long prompt generation
pipe.enable_free_noise(context_length=16, context_stride=4)
pipe.to("cuda")
# Can be a single prompt, or a dictionary with frame timesteps
prompt = {
0: "A caterpillar on a leaf, high quality, photorealistic",
40: "A caterpillar transforming into a cocoon, on a leaf, near flowers, photorealistic",
80: "A cocoon on a leaf, flowers in the backgrond, photorealistic",
120: "A cocoon maturing and a butterfly being born, flowers and leaves visible in the background, photorealistic",
160: "A beautiful butterfly, vibrant colors, sitting on a leaf, flowers in the background, photorealistic",
200: "A beautiful butterfly, flying away in a forest, photorealistic",
240: "A cyberpunk butterfly, neon lights, glowing",
}
negative_prompt = "bad quality, worst quality, jpeg artifacts"
# Run inference
output = pipe(
prompt=prompt,
negative_prompt=negative_prompt,
num_frames=256,
guidance_scale=2.5,
num_inference_steps=10,
generator=torch.Generator("cpu").manual_seed(0),
)
# Save video
frames = output.frames[0]
export_to_video(frames, "output.mp4", fps=16)
```
### FreeNoise memory savings
Since FreeNoise processes multiple frames together, there are parts in the modeling where the memory required exceeds that available on normal consumer GPUs. The main memory bottlenecks that we identified are spatial and temporal attention blocks, upsampling and downsampling blocks, resnet blocks and feed-forward layers. Since most of these blocks operate effectively only on the channel/embedding dimension, one can perform chunked inference across the batch dimensions. The batch dimension in AnimateDiff are either spatial (`[B x F, H x W, C]`) or temporal (`B x H x W, F, C`) in nature (note that it may seem counter-intuitive, but the batch dimension here are correct, because spatial blocks process across the `B x F` dimension while the temporal blocks process across the `B x H x W` dimension). We introduce a `SplitInferenceModule` that makes it easier to chunk across any dimension and perform inference. This saves a lot of memory but comes at the cost of requiring more time for inference.
```diff
# Load pipeline and adapters
# ...
+ pipe.enable_free_noise_split_inference()
+ pipe.unet.enable_forward_chunking(16)
```
The call to `pipe.enable_free_noise_split_inference` method accepts two parameters: `spatial_split_size` (defaults to `256`) and `temporal_split_size` (defaults to `16`). These can be configured based on how much VRAM you have available. A lower split size results in lower memory usage but slower inference, whereas a larger split size results in faster inference at the cost of more memory.
## Using `from_single_file` with the MotionAdapter
@@ -1009,6 +571,7 @@ ckpt_path = "https://huggingface.co/Lightricks/LongAnimateDiff/blob/main/lt_long
adapter = MotionAdapter.from_single_file(ckpt_path, torch_dtype=torch.float16)
pipe = AnimateDiffPipeline.from_pretrained("emilianJR/epiCRealism", motion_adapter=adapter)
```
## AnimateDiffPipeline
@@ -1017,18 +580,6 @@ pipe = AnimateDiffPipeline.from_pretrained("emilianJR/epiCRealism", motion_adapt
- all
- __call__
## AnimateDiffControlNetPipeline
[[autodoc]] AnimateDiffControlNetPipeline
- all
- __call__
## AnimateDiffSparseControlNetPipeline
[[autodoc]] AnimateDiffSparseControlNetPipeline
- all
- __call__
## AnimateDiffSDXLPipeline
[[autodoc]] AnimateDiffSDXLPipeline
@@ -1041,12 +592,6 @@ pipe = AnimateDiffPipeline.from_pretrained("emilianJR/epiCRealism", motion_adapt
- all
- __call__
## AnimateDiffVideoToVideoControlNetPipeline
[[autodoc]] AnimateDiffVideoToVideoControlNetPipeline
- all
- __call__
## AnimateDiffPipelineOutput
[[autodoc]] pipelines.animatediff.AnimateDiffPipelineOutput

View File

@@ -18,7 +18,7 @@ It was developed by the Fal team and more details about it can be found in [this
<Tip>
AuraFlow can be quite expensive to run on consumer hardware devices. However, you can perform a suite of optimizations to run it faster and in a more memory-friendly manner. Check out [this section](https://huggingface.co/blog/sd3#memory-optimizations-for-sd3) for more details.
AuraFlow can be quite expensive to run on consumer hardware devices. However, you can perform a suite of optimizations to run it faster and in a more memory-friendly manner. Check out [this section](https://huggingface.co/blog/sd3#memory-optimizations-for-sd3) for more details.
</Tip>

View File

@@ -1,149 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# CogVideoX
[CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer](https://arxiv.org/abs/2408.06072) from Tsinghua University & ZhipuAI, by Zhuoyi Yang, Jiayan Teng, Wendi Zheng, Ming Ding, Shiyu Huang, Jiazheng Xu, Yuanming Yang, Wenyi Hong, Xiaohan Zhang, Guanyu Feng, Da Yin, Xiaotao Gu, Yuxuan Zhang, Weihan Wang, Yean Cheng, Ting Liu, Bin Xu, Yuxiao Dong, Jie Tang.
The abstract from the paper is:
*We introduce CogVideoX, a large-scale diffusion transformer model designed for generating videos based on text prompts. To efficently model video data, we propose to levearge a 3D Variational Autoencoder (VAE) to compresses videos along both spatial and temporal dimensions. To improve the text-video alignment, we propose an expert transformer with the expert adaptive LayerNorm to facilitate the deep fusion between the two modalities. By employing a progressive training technique, CogVideoX is adept at producing coherent, long-duration videos characterized by significant motion. In addition, we develop an effectively text-video data processing pipeline that includes various data preprocessing strategies and a video captioning method. It significantly helps enhance the performance of CogVideoX, improving both generation quality and semantic alignment. Results show that CogVideoX demonstrates state-of-the-art performance across both multiple machine metrics and human evaluations. The model weight of CogVideoX-2B is publicly available at https://github.com/THUDM/CogVideo.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [zRzRzRzRzRzRzR](https://github.com/zRzRzRzRzRzRzR). The original codebase can be found [here](https://huggingface.co/THUDM). The original weights can be found under [hf.co/THUDM](https://huggingface.co/THUDM).
There are three official CogVideoX checkpoints for text-to-video and video-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`THUDM/CogVideoX-2b`](https://huggingface.co/THUDM/CogVideoX-2b) | torch.float16 |
| [`THUDM/CogVideoX-5b`](https://huggingface.co/THUDM/CogVideoX-5b) | torch.bfloat16 |
| [`THUDM/CogVideoX1.5-5b`](https://huggingface.co/THUDM/CogVideoX1.5-5b) | torch.bfloat16 |
There are two official CogVideoX checkpoints available for image-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`THUDM/CogVideoX-5b-I2V`](https://huggingface.co/THUDM/CogVideoX-5b-I2V) | torch.bfloat16 |
| [`THUDM/CogVideoX-1.5-5b-I2V`](https://huggingface.co/THUDM/CogVideoX-1.5-5b-I2V) | torch.bfloat16 |
For the CogVideoX 1.5 series:
- Text-to-video (T2V) works best at a resolution of 1360x768 because it was trained with that specific resolution.
- Image-to-video (I2V) works for multiple resolutions. The width can vary from 768 to 1360, but the height must be 768. The height/width must be divisible by 16.
- Both T2V and I2V models support generation with 81 and 161 frames and work best at this value. Exporting videos at 16 FPS is recommended.
There are two official CogVideoX checkpoints that support pose controllable generation (by the [Alibaba-PAI](https://huggingface.co/alibaba-pai) team).
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`alibaba-pai/CogVideoX-Fun-V1.1-2b-Pose`](https://huggingface.co/alibaba-pai/CogVideoX-Fun-V1.1-2b-Pose) | torch.bfloat16 |
| [`alibaba-pai/CogVideoX-Fun-V1.1-5b-Pose`](https://huggingface.co/alibaba-pai/CogVideoX-Fun-V1.1-5b-Pose) | torch.bfloat16 |
## Inference
Use [`torch.compile`](https://huggingface.co/docs/diffusers/main/en/tutorials/fast_diffusion#torchcompile) to reduce the inference latency.
First, load the pipeline:
```python
import torch
from diffusers import CogVideoXPipeline, CogVideoXImageToVideoPipeline
from diffusers.utils import export_to_video,load_image
pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b").to("cuda") # or "THUDM/CogVideoX-2b"
```
If you are using the image-to-video pipeline, load it as follows:
```python
pipe = CogVideoXImageToVideoPipeline.from_pretrained("THUDM/CogVideoX-5b-I2V").to("cuda")
```
Then change the memory layout of the pipelines `transformer` component to `torch.channels_last`:
```python
pipe.transformer.to(memory_format=torch.channels_last)
```
Compile the components and run inference:
```python
pipe.transformer = torch.compile(pipeline.transformer, mode="max-autotune", fullgraph=True)
# CogVideoX works well with long and well-described prompts
prompt = "A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical atmosphere of this unique musical performance."
video = pipe(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
```
The [T2V benchmark](https://gist.github.com/a-r-r-o-w/5183d75e452a368fd17448fcc810bd3f) results on an 80GB A100 machine are:
```
Without torch.compile(): Average inference time: 96.89 seconds.
With torch.compile(): Average inference time: 76.27 seconds.
```
### Memory optimization
CogVideoX-2b requires about 19 GB of GPU memory to decode 49 frames (6 seconds of video at 8 FPS) with output resolution 720x480 (W x H), which makes it not possible to run on consumer GPUs or free-tier T4 Colab. The following memory optimizations could be used to reduce the memory footprint. For replication, you can refer to [this](https://gist.github.com/a-r-r-o-w/3959a03f15be5c9bd1fe545b09dfcc93) script.
- `pipe.enable_model_cpu_offload()`:
- Without enabling cpu offloading, memory usage is `33 GB`
- With enabling cpu offloading, memory usage is `19 GB`
- `pipe.enable_sequential_cpu_offload()`:
- Similar to `enable_model_cpu_offload` but can significantly reduce memory usage at the cost of slow inference
- When enabled, memory usage is under `4 GB`
- `pipe.vae.enable_tiling()`:
- With enabling cpu offloading and tiling, memory usage is `11 GB`
- `pipe.vae.enable_slicing()`
### Quantized inference
[torchao](https://github.com/pytorch/ao) and [optimum-quanto](https://github.com/huggingface/optimum-quanto/) can be used to quantize the text encoder, transformer and VAE modules to lower the memory requirements. This makes it possible to run the model on a free-tier T4 Colab or lower VRAM GPUs!
It is also worth noting that torchao quantization is fully compatible with [torch.compile](/optimization/torch2.0#torchcompile), which allows for much faster inference speed. Additionally, models can be serialized and stored in a quantized datatype to save disk space with torchao. Find examples and benchmarks in the gists below.
- [torchao](https://gist.github.com/a-r-r-o-w/4d9732d17412888c885480c6521a9897)
- [quanto](https://gist.github.com/a-r-r-o-w/31be62828b00a9292821b85c1017effa)
## CogVideoXPipeline
[[autodoc]] CogVideoXPipeline
- all
- __call__
## CogVideoXImageToVideoPipeline
[[autodoc]] CogVideoXImageToVideoPipeline
- all
- __call__
## CogVideoXVideoToVideoPipeline
[[autodoc]] CogVideoXVideoToVideoPipeline
- all
- __call__
## CogVideoXFunControlPipeline
[[autodoc]] CogVideoXFunControlPipeline
- all
- __call__
## CogVideoXPipelineOutput
[[autodoc]] pipelines.cogvideo.pipeline_output.CogVideoXPipelineOutput

View File

@@ -1,40 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# CogView3Plus
[CogView3: Finer and Faster Text-to-Image Generation via Relay Diffusion](https://huggingface.co/papers/2403.05121) from Tsinghua University & ZhipuAI, by Wendi Zheng, Jiayan Teng, Zhuoyi Yang, Weihan Wang, Jidong Chen, Xiaotao Gu, Yuxiao Dong, Ming Ding, Jie Tang.
The abstract from the paper is:
*Recent advancements in text-to-image generative systems have been largely driven by diffusion models. However, single-stage text-to-image diffusion models still face challenges, in terms of computational efficiency and the refinement of image details. To tackle the issue, we propose CogView3, an innovative cascaded framework that enhances the performance of text-to-image diffusion. CogView3 is the first model implementing relay diffusion in the realm of text-to-image generation, executing the task by first creating low-resolution images and subsequently applying relay-based super-resolution. This methodology not only results in competitive text-to-image outputs but also greatly reduces both training and inference costs. Our experimental results demonstrate that CogView3 outperforms SDXL, the current state-of-the-art open-source text-to-image diffusion model, by 77.0% in human evaluations, all while requiring only about 1/2 of the inference time. The distilled variant of CogView3 achieves comparable performance while only utilizing 1/10 of the inference time by SDXL.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [zRzRzRzRzRzRzR](https://github.com/zRzRzRzRzRzRzR). The original codebase can be found [here](https://huggingface.co/THUDM). The original weights can be found under [hf.co/THUDM](https://huggingface.co/THUDM).
## CogView3PlusPipeline
[[autodoc]] CogView3PlusPipeline
- all
- __call__
## CogView3PipelineOutput
[[autodoc]] pipelines.cogview3.pipeline_output.CogView3PipelineOutput

View File

@@ -1,56 +0,0 @@
<!--Copyright 2024 The HuggingFace Team, The InstantX Team, and the XLabs Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNet with Flux.1
FluxControlNetPipeline is an implementation of ControlNet for Flux.1.
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
The abstract from the paper is:
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
This controlnet code is implemented by [The InstantX Team](https://huggingface.co/InstantX). You can find pre-trained checkpoints for Flux-ControlNet in the table below:
| ControlNet type | Developer | Link |
| -------- | ---------- | ---- |
| Canny | [The InstantX Team](https://huggingface.co/InstantX) | [Link](https://huggingface.co/InstantX/FLUX.1-dev-Controlnet-Canny) |
| Depth | [The InstantX Team](https://huggingface.co/InstantX) | [Link](https://huggingface.co/Shakker-Labs/FLUX.1-dev-ControlNet-Depth) |
| Union | [The InstantX Team](https://huggingface.co/InstantX) | [Link](https://huggingface.co/InstantX/FLUX.1-dev-Controlnet-Union) |
XLabs ControlNets are also supported, which was contributed by the [XLabs team](https://huggingface.co/XLabs-AI).
| ControlNet type | Developer | Link |
| -------- | ---------- | ---- |
| Canny | [The XLabs Team](https://huggingface.co/XLabs-AI) | [Link](https://huggingface.co/XLabs-AI/flux-controlnet-canny-diffusers) |
| Depth | [The XLabs Team](https://huggingface.co/XLabs-AI) | [Link](https://huggingface.co/XLabs-AI/flux-controlnet-depth-diffusers) |
| HED | [The XLabs Team](https://huggingface.co/XLabs-AI) | [Link](https://huggingface.co/XLabs-AI/flux-controlnet-hed-diffusers) |
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-components-across-pipelines) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## FluxControlNetPipeline
[[autodoc]] FluxControlNetPipeline
- all
- __call__
## FluxPipelineOutput
[[autodoc]] pipelines.flux.pipeline_output.FluxPipelineOutput

View File

@@ -1,4 +1,4 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
<!--Copyright 2023 The HuggingFace Team and The InstantX Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -22,17 +22,7 @@ The abstract from the paper is:
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
This controlnet code is mainly implemented by [The InstantX Team](https://huggingface.co/InstantX). The inpainting-related code was developed by [The Alimama Creative Team](https://huggingface.co/alimama-creative). You can find pre-trained checkpoints for SD3-ControlNet in the table below:
| ControlNet type | Developer | Link |
| -------- | ---------- | ---- |
| Canny | [The InstantX Team](https://huggingface.co/InstantX) | [Link](https://huggingface.co/InstantX/SD3-Controlnet-Canny) |
| Depth | [The InstantX Team](https://huggingface.co/InstantX) | [Link](https://huggingface.co/InstantX/SD3-Controlnet-Depth) |
| Pose | [The InstantX Team](https://huggingface.co/InstantX) | [Link](https://huggingface.co/InstantX/SD3-Controlnet-Pose) |
| Tile | [The InstantX Team](https://huggingface.co/InstantX) | [Link](https://huggingface.co/InstantX/SD3-Controlnet-Tile) |
| Inpainting | [The AlimamaCreative Team](https://huggingface.co/alimama-creative) | [link](https://huggingface.co/alimama-creative/SD3-Controlnet-Inpainting) |
This code is implemented by [The InstantX Team](https://huggingface.co/InstantX). You can find pre-trained checkpoints for SD3-ControlNet on [The InstantX Team](https://huggingface.co/InstantX) Hub profile.
<Tip>
@@ -45,10 +35,5 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
- all
- __call__
## StableDiffusion3ControlNetInpaintingPipeline
[[autodoc]] pipelines.controlnet_sd3.pipeline_stable_diffusion_3_controlnet_inpainting.StableDiffusion3ControlNetInpaintingPipeline
- all
- __call__
## StableDiffusion3PipelineOutput
[[autodoc]] pipelines.stable_diffusion_3.pipeline_output.StableDiffusion3PipelineOutput

View File

@@ -1,35 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNetUnion
ControlNetUnionModel is an implementation of ControlNet for Stable Diffusion XL.
The ControlNet model was introduced in [ControlNetPlus](https://github.com/xinsir6/ControlNetPlus) by xinsir6. It supports multiple conditioning inputs without increasing computation.
*We design a new architecture that can support 10+ control types in condition text-to-image generation and can generate high resolution images visually comparable with midjourney. The network is based on the original ControlNet architecture, we propose two new modules to: 1 Extend the original ControlNet to support different image conditions using the same network parameter. 2 Support multiple conditions input without increasing computation offload, which is especially important for designers who want to edit image in detail, different conditions use the same condition encoder, without adding extra computations or parameters.*
## StableDiffusionXLControlNetUnionPipeline
[[autodoc]] StableDiffusionXLControlNetUnionPipeline
- all
- __call__
## StableDiffusionXLControlNetUnionImg2ImgPipeline
[[autodoc]] StableDiffusionXLControlNetUnionImg2ImgPipeline
- all
- __call__
## StableDiffusionXLControlNetUnionInpaintPipeline
[[autodoc]] StableDiffusionXLControlNetUnionInpaintPipeline
- all
- __call__

View File

@@ -1,406 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Flux
Flux is a series of text-to-image generation models based on diffusion transformers. To know more about Flux, check out the original [blog post](https://blackforestlabs.ai/announcing-black-forest-labs/) by the creators of Flux, Black Forest Labs.
Original model checkpoints for Flux can be found [here](https://huggingface.co/black-forest-labs). Original inference code can be found [here](https://github.com/black-forest-labs/flux).
<Tip>
Flux can be quite expensive to run on consumer hardware devices. However, you can perform a suite of optimizations to run it faster and in a more memory-friendly manner. Check out [this section](https://huggingface.co/blog/sd3#memory-optimizations-for-sd3) for more details. Additionally, Flux can benefit from quantization for memory efficiency with a trade-off in inference latency. Refer to [this blog post](https://huggingface.co/blog/quanto-diffusers) to learn more. For an exhaustive list of resources, check out [this gist](https://gist.github.com/sayakpaul/b664605caf0aa3bf8585ab109dd5ac9c).
</Tip>
Flux comes in the following variants:
| model type | model id |
|:----------:|:--------:|
| Timestep-distilled | [`black-forest-labs/FLUX.1-schnell`](https://huggingface.co/black-forest-labs/FLUX.1-schnell) |
| Guidance-distilled | [`black-forest-labs/FLUX.1-dev`](https://huggingface.co/black-forest-labs/FLUX.1-dev) |
| Fill Inpainting/Outpainting (Guidance-distilled) | [`black-forest-labs/FLUX.1-Fill-dev`](https://huggingface.co/black-forest-labs/FLUX.1-Fill-dev) |
| Canny Control (Guidance-distilled) | [`black-forest-labs/FLUX.1-Canny-dev`](https://huggingface.co/black-forest-labs/FLUX.1-Canny-dev) |
| Depth Control (Guidance-distilled) | [`black-forest-labs/FLUX.1-Depth-dev`](https://huggingface.co/black-forest-labs/FLUX.1-Depth-dev) |
| Canny Control (LoRA) | [`black-forest-labs/FLUX.1-Canny-dev-lora`](https://huggingface.co/black-forest-labs/FLUX.1-Canny-dev-lora) |
| Depth Control (LoRA) | [`black-forest-labs/FLUX.1-Depth-dev-lora`](https://huggingface.co/black-forest-labs/FLUX.1-Depth-dev-lora) |
| Redux (Adapter) | [`black-forest-labs/FLUX.1-Redux-dev`](https://huggingface.co/black-forest-labs/FLUX.1-Redux-dev) |
All checkpoints have different usage which we detail below.
### Timestep-distilled
* `max_sequence_length` cannot be more than 256.
* `guidance_scale` needs to be 0.
* As this is a timestep-distilled model, it benefits from fewer sampling steps.
```python
import torch
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained("black-forest-labs/FLUX.1-schnell", torch_dtype=torch.bfloat16)
pipe.enable_model_cpu_offload()
prompt = "A cat holding a sign that says hello world"
out = pipe(
prompt=prompt,
guidance_scale=0.,
height=768,
width=1360,
num_inference_steps=4,
max_sequence_length=256,
).images[0]
out.save("image.png")
```
### Guidance-distilled
* The guidance-distilled variant takes about 50 sampling steps for good-quality generation.
* It doesn't have any limitations around the `max_sequence_length`.
```python
import torch
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained("black-forest-labs/FLUX.1-dev", torch_dtype=torch.bfloat16)
pipe.enable_model_cpu_offload()
prompt = "a tiny astronaut hatching from an egg on the moon"
out = pipe(
prompt=prompt,
guidance_scale=3.5,
height=768,
width=1360,
num_inference_steps=50,
).images[0]
out.save("image.png")
```
### Fill Inpainting/Outpainting
* Flux Fill pipeline does not require `strength` as an input like regular inpainting pipelines.
* It supports both inpainting and outpainting.
```python
import torch
from diffusers import FluxFillPipeline
from diffusers.utils import load_image
image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/cup.png")
mask = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/cup_mask.png")
repo_id = "black-forest-labs/FLUX.1-Fill-dev"
pipe = FluxFillPipeline.from_pretrained(repo_id, torch_dtype=torch.bfloat16).to("cuda")
image = pipe(
prompt="a white paper cup",
image=image,
mask_image=mask,
height=1632,
width=1232,
max_sequence_length=512,
generator=torch.Generator("cpu").manual_seed(0)
).images[0]
image.save(f"output.png")
```
### Canny Control
**Note:** `black-forest-labs/Flux.1-Canny-dev` is _not_ a [`ControlNetModel`] model. ControlNet models are a separate component from the UNet/Transformer whose residuals are added to the actual underlying model. Canny Control is an alternate architecture that achieves effectively the same results as a ControlNet model would, by using channel-wise concatenation with input control condition and ensuring the transformer learns structure control by following the condition as closely as possible.
```python
# !pip install -U controlnet-aux
import torch
from controlnet_aux import CannyDetector
from diffusers import FluxControlPipeline
from diffusers.utils import load_image
pipe = FluxControlPipeline.from_pretrained("black-forest-labs/FLUX.1-Canny-dev", torch_dtype=torch.bfloat16).to("cuda")
prompt = "A robot made of exotic candies and chocolates of different kinds. The background is filled with confetti and celebratory gifts."
control_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/robot.png")
processor = CannyDetector()
control_image = processor(control_image, low_threshold=50, high_threshold=200, detect_resolution=1024, image_resolution=1024)
image = pipe(
prompt=prompt,
control_image=control_image,
height=1024,
width=1024,
num_inference_steps=50,
guidance_scale=30.0,
).images[0]
image.save("output.png")
```
Canny Control is also possible with a LoRA variant of this condition. The usage is as follows:
```python
# !pip install -U controlnet-aux
import torch
from controlnet_aux import CannyDetector
from diffusers import FluxControlPipeline
from diffusers.utils import load_image
pipe = FluxControlPipeline.from_pretrained("black-forest-labs/FLUX.1-dev", torch_dtype=torch.bfloat16).to("cuda")
pipe.load_lora_weights("black-forest-labs/FLUX.1-Canny-dev-lora")
prompt = "A robot made of exotic candies and chocolates of different kinds. The background is filled with confetti and celebratory gifts."
control_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/robot.png")
processor = CannyDetector()
control_image = processor(control_image, low_threshold=50, high_threshold=200, detect_resolution=1024, image_resolution=1024)
image = pipe(
prompt=prompt,
control_image=control_image,
height=1024,
width=1024,
num_inference_steps=50,
guidance_scale=30.0,
).images[0]
image.save("output.png")
```
### Depth Control
**Note:** `black-forest-labs/Flux.1-Depth-dev` is _not_ a ControlNet model. [`ControlNetModel`] models are a separate component from the UNet/Transformer whose residuals are added to the actual underlying model. Depth Control is an alternate architecture that achieves effectively the same results as a ControlNet model would, by using channel-wise concatenation with input control condition and ensuring the transformer learns structure control by following the condition as closely as possible.
```python
# !pip install git+https://github.com/huggingface/image_gen_aux
import torch
from diffusers import FluxControlPipeline, FluxTransformer2DModel
from diffusers.utils import load_image
from image_gen_aux import DepthPreprocessor
pipe = FluxControlPipeline.from_pretrained("black-forest-labs/FLUX.1-Depth-dev", torch_dtype=torch.bfloat16).to("cuda")
prompt = "A robot made of exotic candies and chocolates of different kinds. The background is filled with confetti and celebratory gifts."
control_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/robot.png")
processor = DepthPreprocessor.from_pretrained("LiheYoung/depth-anything-large-hf")
control_image = processor(control_image)[0].convert("RGB")
image = pipe(
prompt=prompt,
control_image=control_image,
height=1024,
width=1024,
num_inference_steps=30,
guidance_scale=10.0,
generator=torch.Generator().manual_seed(42),
).images[0]
image.save("output.png")
```
Depth Control is also possible with a LoRA variant of this condition. The usage is as follows:
```python
# !pip install git+https://github.com/huggingface/image_gen_aux
import torch
from diffusers import FluxControlPipeline, FluxTransformer2DModel
from diffusers.utils import load_image
from image_gen_aux import DepthPreprocessor
pipe = FluxControlPipeline.from_pretrained("black-forest-labs/FLUX.1-dev", torch_dtype=torch.bfloat16).to("cuda")
pipe.load_lora_weights("black-forest-labs/FLUX.1-Depth-dev-lora")
prompt = "A robot made of exotic candies and chocolates of different kinds. The background is filled with confetti and celebratory gifts."
control_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/robot.png")
processor = DepthPreprocessor.from_pretrained("LiheYoung/depth-anything-large-hf")
control_image = processor(control_image)[0].convert("RGB")
image = pipe(
prompt=prompt,
control_image=control_image,
height=1024,
width=1024,
num_inference_steps=30,
guidance_scale=10.0,
generator=torch.Generator().manual_seed(42),
).images[0]
image.save("output.png")
```
### Redux
* Flux Redux pipeline is an adapter for FLUX.1 base models. It can be used with both flux-dev and flux-schnell, for image-to-image generation.
* You can first use the `FluxPriorReduxPipeline` to get the `prompt_embeds` and `pooled_prompt_embeds`, and then feed them into the `FluxPipeline` for image-to-image generation.
* When use `FluxPriorReduxPipeline` with a base pipeline, you can set `text_encoder=None` and `text_encoder_2=None` in the base pipeline, in order to save VRAM.
```python
import torch
from diffusers import FluxPriorReduxPipeline, FluxPipeline
from diffusers.utils import load_image
device = "cuda"
dtype = torch.bfloat16
repo_redux = "black-forest-labs/FLUX.1-Redux-dev"
repo_base = "black-forest-labs/FLUX.1-dev"
pipe_prior_redux = FluxPriorReduxPipeline.from_pretrained(repo_redux, torch_dtype=dtype).to(device)
pipe = FluxPipeline.from_pretrained(
repo_base,
text_encoder=None,
text_encoder_2=None,
torch_dtype=torch.bfloat16
).to(device)
image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/style_ziggy/img5.png")
pipe_prior_output = pipe_prior_redux(image)
images = pipe(
guidance_scale=2.5,
num_inference_steps=50,
generator=torch.Generator("cpu").manual_seed(0),
**pipe_prior_output,
).images
images[0].save("flux-redux.png")
```
## Running FP16 inference
Flux can generate high-quality images with FP16 (i.e. to accelerate inference on Turing/Volta GPUs) but produces different outputs compared to FP32/BF16. The issue is that some activations in the text encoders have to be clipped when running in FP16, which affects the overall image. Forcing text encoders to run with FP32 inference thus removes this output difference. See [here](https://github.com/huggingface/diffusers/pull/9097#issuecomment-2272292516) for details.
FP16 inference code:
```python
import torch
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained("black-forest-labs/FLUX.1-schnell", torch_dtype=torch.bfloat16) # can replace schnell with dev
# to run on low vram GPUs (i.e. between 4 and 32 GB VRAM)
pipe.enable_sequential_cpu_offload()
pipe.vae.enable_slicing()
pipe.vae.enable_tiling()
pipe.to(torch.float16) # casting here instead of in the pipeline constructor because doing so in the constructor loads all models into CPU memory at once
prompt = "A cat holding a sign that says hello world"
out = pipe(
prompt=prompt,
guidance_scale=0.,
height=768,
width=1360,
num_inference_steps=4,
max_sequence_length=256,
).images[0]
out.save("image.png")
```
## Single File Loading for the `FluxTransformer2DModel`
The `FluxTransformer2DModel` supports loading checkpoints in the original format shipped by Black Forest Labs. This is also useful when trying to load finetunes or quantized versions of the models that have been published by the community.
<Tip>
`FP8` inference can be brittle depending on the GPU type, CUDA version, and `torch` version that you are using. It is recommended that you use the `optimum-quanto` library in order to run FP8 inference on your machine.
</Tip>
The following example demonstrates how to run Flux with less than 16GB of VRAM.
First install `optimum-quanto`
```shell
pip install optimum-quanto
```
Then run the following example
```python
import torch
from diffusers import FluxTransformer2DModel, FluxPipeline
from transformers import T5EncoderModel, CLIPTextModel
from optimum.quanto import freeze, qfloat8, quantize
bfl_repo = "black-forest-labs/FLUX.1-dev"
dtype = torch.bfloat16
transformer = FluxTransformer2DModel.from_single_file("https://huggingface.co/Kijai/flux-fp8/blob/main/flux1-dev-fp8.safetensors", torch_dtype=dtype)
quantize(transformer, weights=qfloat8)
freeze(transformer)
text_encoder_2 = T5EncoderModel.from_pretrained(bfl_repo, subfolder="text_encoder_2", torch_dtype=dtype)
quantize(text_encoder_2, weights=qfloat8)
freeze(text_encoder_2)
pipe = FluxPipeline.from_pretrained(bfl_repo, transformer=None, text_encoder_2=None, torch_dtype=dtype)
pipe.transformer = transformer
pipe.text_encoder_2 = text_encoder_2
pipe.enable_model_cpu_offload()
prompt = "A cat holding a sign that says hello world"
image = pipe(
prompt,
guidance_scale=3.5,
output_type="pil",
num_inference_steps=20,
generator=torch.Generator("cpu").manual_seed(0)
).images[0]
image.save("flux-fp8-dev.png")
```
## FluxPipeline
[[autodoc]] FluxPipeline
- all
- __call__
## FluxImg2ImgPipeline
[[autodoc]] FluxImg2ImgPipeline
- all
- __call__
## FluxInpaintPipeline
[[autodoc]] FluxInpaintPipeline
- all
- __call__
## FluxControlNetInpaintPipeline
[[autodoc]] FluxControlNetInpaintPipeline
- all
- __call__
## FluxControlNetImg2ImgPipeline
[[autodoc]] FluxControlNetImg2ImgPipeline
- all
- __call__
## FluxControlPipeline
[[autodoc]] FluxControlPipeline
- all
- __call__
## FluxControlImg2ImgPipeline
[[autodoc]] FluxControlImg2ImgPipeline
- all
- __call__
## FluxPriorReduxPipeline
[[autodoc]] FluxPriorReduxPipeline
- all
- __call__
## FluxFillPipeline
[[autodoc]] FluxFillPipeline
- all
- __call__

View File

@@ -1,43 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# HunyuanVideo
[HunyuanVideo](https://www.arxiv.org/abs/2412.03603) by Tencent.
*Recent advancements in video generation have significantly impacted daily life for both individuals and industries. However, the leading video generation models remain closed-source, resulting in a notable performance gap between industry capabilities and those available to the public. In this report, we introduce HunyuanVideo, an innovative open-source video foundation model that demonstrates performance in video generation comparable to, or even surpassing, that of leading closed-source models. HunyuanVideo encompasses a comprehensive framework that integrates several key elements, including data curation, advanced architectural design, progressive model scaling and training, and an efficient infrastructure tailored for large-scale model training and inference. As a result, we successfully trained a video generative model with over 13 billion parameters, making it the largest among all open-source models. We conducted extensive experiments and implemented a series of targeted designs to ensure high visual quality, motion dynamics, text-video alignment, and advanced filming techniques. According to evaluations by professionals, HunyuanVideo outperforms previous state-of-the-art models, including Runway Gen-3, Luma 1.6, and three top-performing Chinese video generative models. By releasing the code for the foundation model and its applications, we aim to bridge the gap between closed-source and open-source communities. This initiative will empower individuals within the community to experiment with their ideas, fostering a more dynamic and vibrant video generation ecosystem. The code is publicly available at [this https URL](https://github.com/Tencent/HunyuanVideo).*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
Recommendations for inference:
- Both text encoders should be in `torch.float16`.
- Transformer should be in `torch.bfloat16`.
- VAE should be in `torch.float16`.
- `num_frames` should be of the form `4 * k + 1`, for example `49` or `129`.
- For smaller resolution images, try lower values of `shift` (between `2.0` to `5.0`) in the [Scheduler](https://huggingface.co/docs/diffusers/main/en/api/schedulers/flow_match_euler_discrete#diffusers.FlowMatchEulerDiscreteScheduler.shift). For larger resolution images, try higher values (between `7.0` and `12.0`). The default value is `7.0` for HunyuanVideo.
- For more information about supported resolutions and other details, please refer to the original repository [here](https://github.com/Tencent/HunyuanVideo/).
## HunyuanVideoPipeline
[[autodoc]] HunyuanVideoPipeline
- all
- __call__
## HunyuanVideoPipelineOutput
[[autodoc]] pipelines.hunyuan_video.pipeline_output.HunyuanVideoPipelineOutput

View File

@@ -14,7 +14,7 @@ specific language governing permissions and limitations under the License.
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/kolors/kolors_header_collage.png)
Kolors is a large-scale text-to-image generation model based on latent diffusion, developed by [the Kuaishou Kolors team](https://github.com/Kwai-Kolors/Kolors). Trained on billions of text-image pairs, Kolors exhibits significant advantages over both open-source and closed-source models in visual quality, complex semantic accuracy, and text rendering for both Chinese and English characters. Furthermore, Kolors supports both Chinese and English inputs, demonstrating strong performance in understanding and generating Chinese-specific content. For more details, please refer to this [technical report](https://github.com/Kwai-Kolors/Kolors/blob/master/imgs/Kolors_paper.pdf).
Kolors is a large-scale text-to-image generation model based on latent diffusion, developed by [the Kuaishou Kolors team](kwai-kolors@kuaishou.com). Trained on billions of text-image pairs, Kolors exhibits significant advantages over both open-source and closed-source models in visual quality, complex semantic accuracy, and text rendering for both Chinese and English characters. Furthermore, Kolors supports both Chinese and English inputs, demonstrating strong performance in understanding and generating Chinese-specific content. For more details, please refer to this [technical report](https://github.com/Kwai-Kolors/Kolors/blob/master/imgs/Kolors_paper.pdf).
The abstract from the technical report is:
@@ -41,75 +41,9 @@ image = pipe(
image.save("kolors_sample.png")
```
### IP Adapter
Kolors needs a different IP Adapter to work, and it uses [Openai-CLIP-336](https://huggingface.co/openai/clip-vit-large-patch14-336) as an image encoder.
<Tip>
Using an IP Adapter with Kolors requires more than 24GB of VRAM. To use it, we recommend using [`~DiffusionPipeline.enable_model_cpu_offload`] on consumer GPUs.
</Tip>
<Tip>
While Kolors is integrated in Diffusers, you need to load the image encoder from a revision to use the safetensor files. You can still use the main branch of the original repository if you're comfortable loading pickle checkpoints.
</Tip>
```python
import torch
from transformers import CLIPVisionModelWithProjection
from diffusers import DPMSolverMultistepScheduler, KolorsPipeline
from diffusers.utils import load_image
image_encoder = CLIPVisionModelWithProjection.from_pretrained(
"Kwai-Kolors/Kolors-IP-Adapter-Plus",
subfolder="image_encoder",
low_cpu_mem_usage=True,
torch_dtype=torch.float16,
revision="refs/pr/4",
)
pipe = KolorsPipeline.from_pretrained(
"Kwai-Kolors/Kolors-diffusers", image_encoder=image_encoder, torch_dtype=torch.float16, variant="fp16"
)
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config, use_karras_sigmas=True)
pipe.load_ip_adapter(
"Kwai-Kolors/Kolors-IP-Adapter-Plus",
subfolder="",
weight_name="ip_adapter_plus_general.safetensors",
revision="refs/pr/4",
image_encoder_folder=None,
)
pipe.enable_model_cpu_offload()
ipa_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/kolors/cat_square.png")
image = pipe(
prompt="best quality, high quality",
negative_prompt="",
guidance_scale=6.5,
num_inference_steps=25,
ip_adapter_image=ipa_image,
).images[0]
image.save("kolors_ipa_sample.png")
```
## KolorsPipeline
[[autodoc]] KolorsPipeline
- all
- __call__
## KolorsImg2ImgPipeline
[[autodoc]] KolorsImg2ImgPipeline
- all
- __call__

View File

@@ -1,77 +0,0 @@
<!-- # Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# Latte
![latte text-to-video](https://github.com/Vchitect/Latte/blob/52bc0029899babbd6e9250384c83d8ed2670ff7a/visuals/latte.gif?raw=true)
[Latte: Latent Diffusion Transformer for Video Generation](https://arxiv.org/abs/2401.03048) from Monash University, Shanghai AI Lab, Nanjing University, and Nanyang Technological University.
The abstract from the paper is:
*We propose a novel Latent Diffusion Transformer, namely Latte, for video generation. Latte first extracts spatio-temporal tokens from input videos and then adopts a series of Transformer blocks to model video distribution in the latent space. In order to model a substantial number of tokens extracted from videos, four efficient variants are introduced from the perspective of decomposing the spatial and temporal dimensions of input videos. To improve the quality of generated videos, we determine the best practices of Latte through rigorous experimental analysis, including video clip patch embedding, model variants, timestep-class information injection, temporal positional embedding, and learning strategies. Our comprehensive evaluation demonstrates that Latte achieves state-of-the-art performance across four standard video generation datasets, i.e., FaceForensics, SkyTimelapse, UCF101, and Taichi-HD. In addition, we extend Latte to text-to-video generation (T2V) task, where Latte achieves comparable results compared to recent T2V models. We strongly believe that Latte provides valuable insights for future research on incorporating Transformers into diffusion models for video generation.*
**Highlights**: Latte is a latent diffusion transformer proposed as a backbone for modeling different modalities (trained for text-to-video generation here). It achieves state-of-the-art performance across four standard video benchmarks - [FaceForensics](https://arxiv.org/abs/1803.09179), [SkyTimelapse](https://arxiv.org/abs/1709.07592), [UCF101](https://arxiv.org/abs/1212.0402) and [Taichi-HD](https://arxiv.org/abs/2003.00196). To prepare and download the datasets for evaluation, please refer to [this https URL](https://github.com/Vchitect/Latte/blob/main/docs/datasets_evaluation.md).
This pipeline was contributed by [maxin-cn](https://github.com/maxin-cn). The original codebase can be found [here](https://github.com/Vchitect/Latte). The original weights can be found under [hf.co/maxin-cn](https://huggingface.co/maxin-cn).
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
### Inference
Use [`torch.compile`](https://huggingface.co/docs/diffusers/main/en/tutorials/fast_diffusion#torchcompile) to reduce the inference latency.
First, load the pipeline:
```python
import torch
from diffusers import LattePipeline
pipeline = LattePipeline.from_pretrained(
"maxin-cn/Latte-1", torch_dtype=torch.float16
).to("cuda")
```
Then change the memory layout of the pipelines `transformer` and `vae` components to `torch.channels-last`:
```python
pipeline.transformer.to(memory_format=torch.channels_last)
pipeline.vae.to(memory_format=torch.channels_last)
```
Finally, compile the components and run inference:
```python
pipeline.transformer = torch.compile(pipeline.transformer)
pipeline.vae.decode = torch.compile(pipeline.vae.decode)
video = pipeline(prompt="A dog wearing sunglasses floating in space, surreal, nebulae in background").frames[0]
```
The [benchmark](https://gist.github.com/a-r-r-o-w/4e1694ca46374793c0361d740a99ff19) results on an 80GB A100 machine are:
```
Without torch.compile(): Average inference time: 16.246 seconds.
With torch.compile(): Average inference time: 14.573 seconds.
```
## LattePipeline
[[autodoc]] LattePipeline
- all
- __call__

View File

@@ -1,80 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# LTX
[LTX Video](https://huggingface.co/Lightricks/LTX-Video) is the first DiT-based video generation model capable of generating high-quality videos in real-time. It produces 24 FPS videos at a 768x512 resolution faster than they can be watched. Trained on a large-scale dataset of diverse videos, the model generates high-resolution videos with realistic and varied content. We provide a model for both text-to-video as well as image + text-to-video usecases.
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Loading Single Files
Loading the original LTX Video checkpoints is also possible with [`~ModelMixin.from_single_file`].
```python
import torch
from diffusers import AutoencoderKLLTXVideo, LTXImageToVideoPipeline, LTXVideoTransformer3DModel
single_file_url = "https://huggingface.co/Lightricks/LTX-Video/ltx-video-2b-v0.9.safetensors"
transformer = LTXVideoTransformer3DModel.from_single_file(
single_file_url, torch_dtype=torch.bfloat16
)
vae = AutoencoderKLLTXVideo.from_single_file(single_file_url, torch_dtype=torch.bfloat16)
pipe = LTXImageToVideoPipeline.from_pretrained(
"Lightricks/LTX-Video", transformer=transformer, vae=vae, torch_dtype=torch.bfloat16
)
# ... inference code ...
```
Alternatively, the pipeline can be used to load the weights with [`~FromSingleFileMixin.from_single_file`].
```python
import torch
from diffusers import LTXImageToVideoPipeline
from transformers import T5EncoderModel, T5Tokenizer
single_file_url = "https://huggingface.co/Lightricks/LTX-Video/ltx-video-2b-v0.9.safetensors"
text_encoder = T5EncoderModel.from_pretrained(
"Lightricks/LTX-Video", subfolder="text_encoder", torch_dtype=torch.bfloat16
)
tokenizer = T5Tokenizer.from_pretrained(
"Lightricks/LTX-Video", subfolder="tokenizer", torch_dtype=torch.bfloat16
)
pipe = LTXImageToVideoPipeline.from_single_file(
single_file_url, text_encoder=text_encoder, tokenizer=tokenizer, torch_dtype=torch.bfloat16
)
```
Refer to [this section](https://huggingface.co/docs/diffusers/main/en/api/pipelines/cogvideox#memory-optimization) to learn more about optimizing memory consumption.
## LTXPipeline
[[autodoc]] LTXPipeline
- all
- __call__
## LTXImageToVideoPipeline
[[autodoc]] LTXImageToVideoPipeline
- all
- __call__
## LTXPipelineOutput
[[autodoc]] pipelines.ltx.pipeline_output.LTXPipelineOutput

View File

@@ -43,8 +43,6 @@ Lumina-T2X has the following components:
* It uses a Flow-based Large Diffusion Transformer as the backbone
* It supports different any modalities with one backbone and corresponding encoder, decoder.
This pipeline was contributed by [PommesPeter](https://github.com/PommesPeter). The original codebase can be found [here](https://github.com/Alpha-VLLM/Lumina-T2X). The original weights can be found under [hf.co/Alpha-VLLM](https://huggingface.co/Alpha-VLLM).
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
@@ -59,7 +57,7 @@ First, load the pipeline:
```python
from diffusers import LuminaText2ImgPipeline
import torch
import torch
pipeline = LuminaText2ImgPipeline.from_pretrained(
"Alpha-VLLM/Lumina-Next-SFT-diffusers", torch_dtype=torch.bfloat16
@@ -87,4 +85,4 @@ image = pipeline(prompt="Upper body of a young woman in a Victorian-era outfit w
[[autodoc]] LuminaText2ImgPipeline
- all
- __call__

View File

@@ -1,36 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# Mochi
[Mochi 1 Preview](https://huggingface.co/genmo/mochi-1-preview) from Genmo.
*Mochi 1 preview is an open state-of-the-art video generation model with high-fidelity motion and strong prompt adherence in preliminary evaluation. This model dramatically closes the gap between closed and open video generation systems. The model is released under a permissive Apache 2.0 license.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## MochiPipeline
[[autodoc]] MochiPipeline
- all
- __call__
## MochiPipelineOutput
[[autodoc]] pipelines.mochi.pipeline_output.MochiPipelineOutput

View File

@@ -30,64 +30,62 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| Pipeline | Tasks |
|---|---|
| [aMUSEd](amused) | text2image |
| [AltDiffusion](alt_diffusion) | image2image |
| [AnimateDiff](animatediff) | text2video |
| [Attend-and-Excite](attend_and_excite) | text2image |
| [Audio Diffusion](audio_diffusion) | image2audio |
| [AudioLDM](audioldm) | text2audio |
| [AudioLDM2](audioldm2) | text2audio |
| [AuraFlow](auraflow) | text2image |
| [BLIP Diffusion](blip_diffusion) | text2image |
| [CogVideoX](cogvideox) | text2video |
| [Consistency Models](consistency_models) | unconditional image generation |
| [ControlNet](controlnet) | text2image, image2image, inpainting |
| [ControlNet with Flux.1](controlnet_flux) | text2image |
| [ControlNet with Hunyuan-DiT](controlnet_hunyuandit) | text2image |
| [ControlNet with Stable Diffusion 3](controlnet_sd3) | text2image |
| [ControlNet with Stable Diffusion XL](controlnet_sdxl) | text2image |
| [ControlNet-XS](controlnetxs) | text2image |
| [ControlNet-XS with Stable Diffusion XL](controlnetxs_sdxl) | text2image |
| [Cycle Diffusion](cycle_diffusion) | image2image |
| [Dance Diffusion](dance_diffusion) | unconditional audio generation |
| [DDIM](ddim) | unconditional image generation |
| [DDPM](ddpm) | unconditional image generation |
| [DeepFloyd IF](deepfloyd_if) | text2image, image2image, inpainting, super-resolution |
| [DiffEdit](diffedit) | inpainting |
| [DiT](dit) | text2image |
| [Flux](flux) | text2image |
| [Hunyuan-DiT](hunyuandit) | text2image |
| [I2VGen-XL](i2vgenxl) | text2video |
| [GLIGEN](stable_diffusion/gligen) | text2image |
| [InstructPix2Pix](pix2pix) | image editing |
| [Kandinsky 2.1](kandinsky) | text2image, image2image, inpainting, interpolation |
| [Kandinsky 2.2](kandinsky_v22) | text2image, image2image, inpainting |
| [Kandinsky 3](kandinsky3) | text2image, image2image |
| [Kolors](kolors) | text2image |
| [Latent Consistency Models](latent_consistency_models) | text2image |
| [Latent Diffusion](latent_diffusion) | text2image, super-resolution |
| [Latte](latte) | text2image |
| [LDM3D](stable_diffusion/ldm3d_diffusion) | text2image, text-to-3D, text-to-pano, upscaling |
| [LEDITS++](ledits_pp) | image editing |
| [Lumina-T2X](lumina) | text2image |
| [Marigold](marigold) | depth |
| [MultiDiffusion](panorama) | text2image |
| [MusicLDM](musicldm) | text2audio |
| [PAG](pag) | text2image |
| [Paint by Example](paint_by_example) | inpainting |
| [PIA](pia) | image2video |
| [ParaDiGMS](paradigms) | text2image |
| [Pix2Pix Zero](pix2pix_zero) | image editing |
| [PixArt-α](pixart) | text2image |
| [PixArt-Σ](pixart_sigma) | text2image |
| [PNDM](pndm) | unconditional image generation |
| [RePaint](repaint) | inpainting |
| [Score SDE VE](score_sde_ve) | unconditional image generation |
| [Self-Attention Guidance](self_attention_guidance) | text2image |
| [Semantic Guidance](semantic_stable_diffusion) | text2image |
| [Shap-E](shap_e) | text-to-3D, image-to-3D |
| [Stable Audio](stable_audio) | text2audio |
| [Stable Cascade](stable_cascade) | text2image |
| [Spectrogram Diffusion](spectrogram_diffusion) | |
| [Stable Diffusion](stable_diffusion/overview) | text2image, image2image, depth2image, inpainting, image variation, latent upscaler, super-resolution |
| [Stable Diffusion Model Editing](model_editing) | model editing |
| [Stable Diffusion XL](stable_diffusion/stable_diffusion_xl) | text2image, image2image, inpainting |
| [Stable Diffusion XL Turbo](stable_diffusion/sdxl_turbo) | text2image, image2image, inpainting |
| [Stable unCLIP](stable_unclip) | text2image, image variation |
| [Stochastic Karras VE](stochastic_karras_ve) | unconditional image generation |
| [T2I-Adapter](stable_diffusion/adapter) | text2image |
| [Text2Video](text_to_video) | text2video, video2video |
| [Text2Video-Zero](text_to_video_zero) | text2video |
| [unCLIP](unclip) | text2image, image variation |
| [Unconditional Latent Diffusion](latent_diffusion_uncond) | unconditional image generation |
| [UniDiffuser](unidiffuser) | text2image, image2text, image variation, text variation, unconditional image generation, unconditional audio generation |
| [Value-guided planning](value_guided_sampling) | value guided sampling |
| [Versatile Diffusion](versatile_diffusion) | text2image, image variation |
| [VQ Diffusion](vq_diffusion) | text2image |
| [Wuerstchen](wuerstchen) | text2image |
## DiffusionPipeline

View File

@@ -20,54 +20,13 @@ The abstract from the paper is:
*Recent studies have demonstrated that diffusion models are capable of generating high-quality samples, but their quality heavily depends on sampling guidance techniques, such as classifier guidance (CG) and classifier-free guidance (CFG). These techniques are often not applicable in unconditional generation or in various downstream tasks such as image restoration. In this paper, we propose a novel sampling guidance, called Perturbed-Attention Guidance (PAG), which improves diffusion sample quality across both unconditional and conditional settings, achieving this without requiring additional training or the integration of external modules. PAG is designed to progressively enhance the structure of samples throughout the denoising process. It involves generating intermediate samples with degraded structure by substituting selected self-attention maps in diffusion U-Net with an identity matrix, by considering the self-attention mechanisms' ability to capture structural information, and guiding the denoising process away from these degraded samples. In both ADM and Stable Diffusion, PAG surprisingly improves sample quality in conditional and even unconditional scenarios. Moreover, PAG significantly improves the baseline performance in various downstream tasks where existing guidances such as CG or CFG cannot be fully utilized, including ControlNet with empty prompts and image restoration such as inpainting and deblurring.*
PAG can be used by specifying the `pag_applied_layers` as a parameter when instantiating a PAG pipeline. It can be a single string or a list of strings. Each string can be a unique layer identifier or a regular expression to identify one or more layers.
- Full identifier as a normal string: `down_blocks.2.attentions.0.transformer_blocks.0.attn1.processor`
- Full identifier as a RegEx: `down_blocks.2.(attentions|motion_modules).0.transformer_blocks.0.attn1.processor`
- Partial identifier as a RegEx: `down_blocks.2`, or `attn1`
- List of identifiers (can be combo of strings and ReGex): `["blocks.1", "blocks.(14|20)", r"down_blocks\.(2,3)"]`
<Tip warning={true}>
Since RegEx is supported as a way for matching layer identifiers, it is crucial to use it correctly otherwise there might be unexpected behaviour. The recommended way to use PAG is by specifying layers as `blocks.{layer_index}` and `blocks.({layer_index_1|layer_index_2|...})`. Using it in any other way, while doable, may bypass our basic validation checks and give you unexpected results.
</Tip>
## AnimateDiffPAGPipeline
[[autodoc]] AnimateDiffPAGPipeline
- all
- __call__
## HunyuanDiTPAGPipeline
[[autodoc]] HunyuanDiTPAGPipeline
- all
- __call__
## KolorsPAGPipeline
[[autodoc]] KolorsPAGPipeline
- all
- __call__
## StableDiffusionPAGInpaintPipeline
[[autodoc]] StableDiffusionPAGInpaintPipeline
- all
- __call__
## StableDiffusionPAGPipeline
[[autodoc]] StableDiffusionPAGPipeline
- all
- __call__
## StableDiffusionPAGImg2ImgPipeline
[[autodoc]] StableDiffusionPAGImg2ImgPipeline
- all
- __call__
## StableDiffusionControlNetPAGPipeline
[[autodoc]] StableDiffusionControlNetPAGPipeline
## StableDiffusionControlNetPAGInpaintPipeline
[[autodoc]] StableDiffusionControlNetPAGInpaintPipeline
- all
- __call__
@@ -90,23 +49,3 @@ Since RegEx is supported as a way for matching layer identifiers, it is crucial
[[autodoc]] StableDiffusionXLControlNetPAGPipeline
- all
- __call__
## StableDiffusionXLControlNetPAGImg2ImgPipeline
[[autodoc]] StableDiffusionXLControlNetPAGImg2ImgPipeline
- all
- __call__
## StableDiffusion3PAGPipeline
[[autodoc]] StableDiffusion3PAGPipeline
- all
- __call__
## StableDiffusion3PAGImg2ImgPipeline
[[autodoc]] StableDiffusion3PAGImg2ImgPipeline
- all
- __call__
## PixArtSigmaPAGPipeline
[[autodoc]] PixArtSigmaPAGPipeline
- all
- __call__

View File

@@ -1,67 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# SanaPipeline
[SANA: Efficient High-Resolution Image Synthesis with Linear Diffusion Transformers](https://huggingface.co/papers/2410.10629) from NVIDIA and MIT HAN Lab, by Enze Xie, Junsong Chen, Junyu Chen, Han Cai, Haotian Tang, Yujun Lin, Zhekai Zhang, Muyang Li, Ligeng Zhu, Yao Lu, Song Han.
The abstract from the paper is:
*We introduce Sana, a text-to-image framework that can efficiently generate images up to 4096×4096 resolution. Sana can synthesize high-resolution, high-quality images with strong text-image alignment at a remarkably fast speed, deployable on laptop GPU. Core designs include: (1) Deep compression autoencoder: unlike traditional AEs, which compress images only 8×, we trained an AE that can compress images 32×, effectively reducing the number of latent tokens. (2) Linear DiT: we replace all vanilla attention in DiT with linear attention, which is more efficient at high resolutions without sacrificing quality. (3) Decoder-only text encoder: we replaced T5 with modern decoder-only small LLM as the text encoder and designed complex human instruction with in-context learning to enhance the image-text alignment. (4) Efficient training and sampling: we propose Flow-DPM-Solver to reduce sampling steps, with efficient caption labeling and selection to accelerate convergence. As a result, Sana-0.6B is very competitive with modern giant diffusion model (e.g. Flux-12B), being 20 times smaller and 100+ times faster in measured throughput. Moreover, Sana-0.6B can be deployed on a 16GB laptop GPU, taking less than 1 second to generate a 1024×1024 resolution image. Sana enables content creation at low cost. Code and model will be publicly released.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [lawrence-cj](https://github.com/lawrence-cj) and [chenjy2003](https://github.com/chenjy2003). The original codebase can be found [here](https://github.com/NVlabs/Sana). The original weights can be found under [hf.co/Efficient-Large-Model](https://huggingface.co/Efficient-Large-Model).
Available models:
| Model | Recommended dtype |
|:-----:|:-----------------:|
| [`Efficient-Large-Model/Sana_1600M_1024px_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_1600M_1024px_diffusers) | `torch.float16` |
| [`Efficient-Large-Model/Sana_1600M_1024px_MultiLing_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_1600M_1024px_MultiLing_diffusers) | `torch.float16` |
| [`Efficient-Large-Model/Sana_1600M_1024px_BF16_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_1600M_1024px_BF16_diffusers) | `torch.bfloat16` |
| [`Efficient-Large-Model/Sana_1600M_512px_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_1600M_512px_diffusers) | `torch.float16` |
| [`Efficient-Large-Model/Sana_1600M_512px_MultiLing_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_1600M_512px_MultiLing_diffusers) | `torch.float16` |
| [`Efficient-Large-Model/Sana_600M_1024px_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_600M_1024px_diffusers) | `torch.float16` |
| [`Efficient-Large-Model/Sana_600M_512px_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_600M_512px_diffusers) | `torch.float16` |
Refer to [this](https://huggingface.co/collections/Efficient-Large-Model/sana-673efba2a57ed99843f11f9e) collection for more information.
Note: The recommended dtype mentioned is for the transformer weights. The text encoder and VAE weights must stay in `torch.bfloat16` or `torch.float32` for the model to work correctly. Please refer to the inference example below to see how to load the model with the recommended dtype.
<Tip>
Make sure to pass the `variant` argument for downloaded checkpoints to use lower disk space. Set it to `"fp16"` for models with recommended dtype as `torch.float16`, and `"bf16"` for models with recommended dtype as `torch.bfloat16`. By default, `torch.float32` weights are downloaded, which use twice the amount of disk storage. Additionally, `torch.float32` weights can be downcasted on-the-fly by specifying the `torch_dtype` argument. Read about it in the [docs](https://huggingface.co/docs/diffusers/v0.31.0/en/api/pipelines/overview#diffusers.DiffusionPipeline.from_pretrained).
</Tip>
## SanaPipeline
[[autodoc]] SanaPipeline
- all
- __call__
## SanaPAGPipeline
[[autodoc]] SanaPAGPipeline
- all
- __call__
## SanaPipelineOutput
[[autodoc]] pipelines.sana.pipeline_output.SanaPipelineOutput

View File

@@ -1,42 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Stable Audio
Stable Audio was proposed in [Stable Audio Open](https://arxiv.org/abs/2407.14358) by Zach Evans et al. . it takes a text prompt as input and predicts the corresponding sound or music sample.
Stable Audio Open generates variable-length (up to 47s) stereo audio at 44.1kHz from text prompts. It comprises three components: an autoencoder that compresses waveforms into a manageable sequence length, a T5-based text embedding for text conditioning, and a transformer-based diffusion (DiT) model that operates in the latent space of the autoencoder.
Stable Audio is trained on a corpus of around 48k audio recordings, where around 47k are from Freesound and the rest are from the Free Music Archive (FMA). All audio files are licensed under CC0, CC BY, or CC Sampling+. This data is used to train the autoencoder and the DiT.
The abstract of the paper is the following:
*Open generative models are vitally important for the community, allowing for fine-tunes and serving as baselines when presenting new models. However, most current text-to-audio models are private and not accessible for artists and researchers to build upon. Here we describe the architecture and training process of a new open-weights text-to-audio model trained with Creative Commons data. Our evaluation shows that the model's performance is competitive with the state-of-the-art across various metrics. Notably, the reported FDopenl3 results (measuring the realism of the generations) showcase its potential for high-quality stereo sound synthesis at 44.1kHz.*
This pipeline was contributed by [Yoach Lacombe](https://huggingface.co/ylacombe). The original codebase can be found at [Stability-AI/stable-audio-tools](https://github.com/Stability-AI/stable-audio-tools).
## Tips
When constructing a prompt, keep in mind:
* Descriptive prompt inputs work best; use adjectives to describe the sound (for example, "high quality" or "clear") and make the prompt context specific where possible (e.g. "melodic techno with a fast beat and synths" works better than "techno").
* Using a *negative prompt* can significantly improve the quality of the generated audio. Try using a negative prompt of "low quality, average quality".
During inference:
* The _quality_ of the generated audio sample can be controlled by the `num_inference_steps` argument; higher steps give higher quality audio at the expense of slower inference.
* Multiple waveforms can be generated in one go: set `num_waveforms_per_prompt` to a value greater than 1 to enable. Automatic scoring will be performed between the generated waveforms and prompt text, and the audios ranked from best to worst accordingly.
## StableAudioPipeline
[[autodoc]] StableAudioPipeline
- all
- __call__

View File

@@ -19,7 +19,7 @@ The Stable Diffusion model can also be applied to inpainting which lets you edit
It is recommended to use this pipeline with checkpoints that have been specifically fine-tuned for inpainting, such
as [runwayml/stable-diffusion-inpainting](https://huggingface.co/runwayml/stable-diffusion-inpainting). Default
text-to-image Stable Diffusion checkpoints, such as
[stable-diffusion-v1-5/stable-diffusion-v1-5](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) are also compatible but they might be less performant.
[runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5) are also compatible but they might be less performant.
<Tip>

View File

@@ -203,7 +203,7 @@ from diffusers import StableDiffusionImg2ImgPipeline
import gradio as gr
pipe = StableDiffusionImg2ImgPipeline.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5")
pipe = StableDiffusionImg2ImgPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
gr.Interface.from_pipeline(pipe).launch()
```

View File

@@ -54,11 +54,6 @@ image = pipe(
image.save("sd3_hello_world.png")
```
**Note:** Stable Diffusion 3.5 can also be run using the SD3 pipeline, and all mentioned optimizations and techniques apply to it as well. In total there are three official models in the SD3 family:
- [`stabilityai/stable-diffusion-3-medium-diffusers`](https://huggingface.co/stabilityai/stable-diffusion-3-medium-diffusers)
- [`stabilityai/stable-diffusion-3.5-large`](https://huggingface.co/stabilityai/stable-diffusion-3-5-large)
- [`stabilityai/stable-diffusion-3.5-large-turbo`](https://huggingface.co/stabilityai/stable-diffusion-3-5-large-turbo)
## Memory Optimisations for SD3
SD3 uses three text encoders, one if which is the very large T5-XXL model. This makes it challenging to run the model on GPUs with less than 24GB of VRAM, even when using `fp16` precision. The following section outlines a few memory optimizations in Diffusers that make it easier to run SD3 on low resource hardware.
@@ -313,26 +308,6 @@ image = pipe("a picture of a cat holding a sign that says hello world").images[0
image.save('sd3-single-file-t5-fp8.png')
```
### Loading the single file checkpoint for the Stable Diffusion 3.5 Transformer Model
```python
import torch
from diffusers import SD3Transformer2DModel, StableDiffusion3Pipeline
transformer = SD3Transformer2DModel.from_single_file(
"https://huggingface.co/stabilityai/stable-diffusion-3.5-large-turbo/blob/main/sd3.5_large.safetensors",
torch_dtype=torch.bfloat16,
)
pipe = StableDiffusion3Pipeline.from_pretrained(
"stabilityai/stable-diffusion-3.5-large",
transformer=transformer,
torch_dtype=torch.bfloat16,
)
pipe.enable_model_cpu_offload()
image = pipe("a cat holding a sign that says hello world").images[0]
image.save("sd35.png")
```
## StableDiffusion3Pipeline
[[autodoc]] StableDiffusion3Pipeline

View File

@@ -40,9 +40,8 @@ To generate a video from prompt, run the following Python code:
```python
import torch
from diffusers import TextToVideoZeroPipeline
import imageio
model_id = "stable-diffusion-v1-5/stable-diffusion-v1-5"
model_id = "runwayml/stable-diffusion-v1-5"
pipe = TextToVideoZeroPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
prompt = "A panda is playing guitar on times square"
@@ -64,7 +63,7 @@ import torch
from diffusers import TextToVideoZeroPipeline
import numpy as np
model_id = "stable-diffusion-v1-5/stable-diffusion-v1-5"
model_id = "runwayml/stable-diffusion-v1-5"
pipe = TextToVideoZeroPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
seed = 0
video_length = 24 #24 ÷ 4fps = 6 seconds
@@ -138,7 +137,7 @@ To generate a video from prompt with additional pose control
from diffusers import StableDiffusionControlNetPipeline, ControlNetModel
from diffusers.pipelines.text_to_video_synthesis.pipeline_text_to_video_zero import CrossFrameAttnProcessor
model_id = "stable-diffusion-v1-5/stable-diffusion-v1-5"
model_id = "runwayml/stable-diffusion-v1-5"
controlnet = ControlNetModel.from_pretrained("lllyasviel/sd-controlnet-openpose", torch_dtype=torch.float16)
pipe = StableDiffusionControlNetPipeline.from_pretrained(
model_id, controlnet=controlnet, torch_dtype=torch.float16

View File

@@ -1,40 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Quantization
Quantization techniques reduce memory and computational costs by representing weights and activations with lower-precision data types like 8-bit integers (int8). This enables loading larger models you normally wouldn't be able to fit into memory, and speeding up inference. Diffusers supports 8-bit and 4-bit quantization with [bitsandbytes](https://huggingface.co/docs/bitsandbytes/en/index).
Quantization techniques that aren't supported in Transformers can be added with the [`DiffusersQuantizer`] class.
<Tip>
Learn how to quantize models in the [Quantization](../quantization/overview) guide.
</Tip>
## BitsAndBytesConfig
[[autodoc]] BitsAndBytesConfig
## GGUFQuantizationConfig
[[autodoc]] GGUFQuantizationConfig
## TorchAoConfig
[[autodoc]] TorchAoConfig
## DiffusersQuantizer
[[autodoc]] quantizers.base.DiffusersQuantizer

View File

@@ -1,24 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# CosineDPMSolverMultistepScheduler
The [`CosineDPMSolverMultistepScheduler`] is a variant of [`DPMSolverMultistepScheduler`] with cosine schedule, proposed by Nichol and Dhariwal (2021).
It is being used in the [Stable Audio Open](https://arxiv.org/abs/2407.14358) paper and the [Stability-AI/stable-audio-tool](https://github.com/Stability-AI/stable-audio-tool) codebase.
This scheduler was contributed by [Yoach Lacombe](https://huggingface.co/ylacombe).
## CosineDPMSolverMultistepScheduler
[[autodoc]] CosineDPMSolverMultistepScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput

View File

@@ -45,15 +45,6 @@ Many schedulers are implemented from the [k-diffusion](https://github.com/crowso
| N/A | [`DEISMultistepScheduler`] | |
| N/A | [`UniPCMultistepScheduler`] | |
## Noise schedules and schedule types
| A1111/k-diffusion | 🤗 Diffusers |
|--------------------------|----------------------------------------------------------------------------|
| Karras | init with `use_karras_sigmas=True` |
| sgm_uniform | init with `timestep_spacing="trailing"` |
| simple | init with `timestep_spacing="trailing"` |
| exponential | init with `timestep_spacing="linspace"`, `use_exponential_sigmas=True` |
| beta | init with `timestep_spacing="linspace"`, `use_beta_sigmas=True` |
All schedulers are built from the base [`SchedulerMixin`] class which implements low level utilities shared by all schedulers.
## SchedulerMixin

View File

@@ -1,82 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Community Projects
Welcome to Community Projects. This space is dedicated to showcasing the incredible work and innovative applications created by our vibrant community using the `diffusers` library.
This section aims to:
- Highlight diverse and inspiring projects built with `diffusers`
- Foster knowledge sharing within our community
- Provide real-world examples of how `diffusers` can be leveraged
Happy exploring, and thank you for being part of the Diffusers community!
<table>
<tr>
<th>Project Name</th>
<th>Description</th>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/carson-katri/dream-textures"> dream-textures </a></td>
<td>Stable Diffusion built-in to Blender</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/megvii-research/HiDiffusion"> HiDiffusion </a></td>
<td>Increases the resolution and speed of your diffusion model by only adding a single line of code</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/lllyasviel/IC-Light"> IC-Light </a></td>
<td>IC-Light is a project to manipulate the illumination of images</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/InstantID/InstantID"> InstantID </a></td>
<td>InstantID : Zero-shot Identity-Preserving Generation in Seconds</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/Sanster/IOPaint"> IOPaint </a></td>
<td>Image inpainting tool powered by SOTA AI Model. Remove any unwanted object, defect, people from your pictures or erase and replace(powered by stable diffusion) any thing on your pictures.</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/bmaltais/kohya_ss"> Kohya </a></td>
<td>Gradio GUI for Kohya's Stable Diffusion trainers</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/magic-research/magic-animate"> MagicAnimate </a></td>
<td>MagicAnimate: Temporally Consistent Human Image Animation using Diffusion Model</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/levihsu/OOTDiffusion"> OOTDiffusion </a></td>
<td>Outfitting Fusion based Latent Diffusion for Controllable Virtual Try-on</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/vladmandic/automatic"> SD.Next </a></td>
<td>SD.Next: Advanced Implementation of Stable Diffusion and other Diffusion-based generative image models</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/ashawkey/stable-dreamfusion"> stable-dreamfusion </a></td>
<td>Text-to-3D & Image-to-3D & Mesh Exportation with NeRF + Diffusion</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/HVision-NKU/StoryDiffusion"> StoryDiffusion </a></td>
<td>StoryDiffusion can create a magic story by generating consistent images and videos.</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/cumulo-autumn/StreamDiffusion"> StreamDiffusion </a></td>
<td>A Pipeline-Level Solution for Real-Time Interactive Generation</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/Netwrck/stable-diffusion-server"> Stable Diffusion Server </a></td>
<td>A server configured for Inpainting/Generation/img2img with one stable diffusion model</td>
</tr>
</table>

View File

@@ -92,7 +92,7 @@ images = sd_pipeline(sample_prompts, num_images_per_prompt=1, generator=generato
![parti-prompts-14](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/evaluation_diffusion_models/parti-prompts-14.png)
We can also set `num_images_per_prompt` accordingly to compare different images for the same prompt. Running the same pipeline but with a different checkpoint ([v1-5](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5)), yields:
We can also set `num_images_per_prompt` accordingly to compare different images for the same prompt. Running the same pipeline but with a different checkpoint ([v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5)), yields:
![parti-prompts-15](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/evaluation_diffusion_models/parti-prompts-15.png)
@@ -177,11 +177,11 @@ generator = torch.manual_seed(seed)
images = sd_pipeline(prompts, num_images_per_prompt=1, generator=generator, output_type="np").images
```
Then we load the [v1-5 checkpoint](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) to generate images:
Then we load the [v1-5 checkpoint](https://huggingface.co/runwayml/stable-diffusion-v1-5) to generate images:
```python
model_ckpt_1_5 = "stable-diffusion-v1-5/stable-diffusion-v1-5"
sd_pipeline_1_5 = StableDiffusionPipeline.from_pretrained(model_ckpt_1_5, torch_dtype=torch.float16).to("cuda")
model_ckpt_1_5 = "runwayml/stable-diffusion-v1-5"
sd_pipeline_1_5 = StableDiffusionPipeline.from_pretrained(model_ckpt_1_5, torch_dtype=weight_dtype).to(device)
images_1_5 = sd_pipeline_1_5(prompts, num_images_per_prompt=1, generator=generator, output_type="np").images
```
@@ -198,7 +198,7 @@ print(f"CLIP Score with v-1-5: {sd_clip_score_1_5}")
# CLIP Score with v-1-5: 36.2137
```
It seems like the [v1-5](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5) checkpoint performs better than its predecessor. Note, however, that the number of prompts we used to compute the CLIP scores is quite low. For a more practical evaluation, this number should be way higher, and the prompts should be diverse.
It seems like the [v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5) checkpoint performs better than its predecessor. Note, however, that the number of prompts we used to compute the CLIP scores is quite low. For a more practical evaluation, this number should be way higher, and the prompts should be diverse.
<Tip warning={true}>
@@ -280,7 +280,7 @@ from diffusers import StableDiffusionInstructPix2PixPipeline
instruct_pix2pix_pipeline = StableDiffusionInstructPix2PixPipeline.from_pretrained(
"timbrooks/instruct-pix2pix", torch_dtype=torch.float16
).to("cuda")
).to(device)
```
Now, we perform the edits:
@@ -326,9 +326,9 @@ from transformers import (
clip_id = "openai/clip-vit-large-patch14"
tokenizer = CLIPTokenizer.from_pretrained(clip_id)
text_encoder = CLIPTextModelWithProjection.from_pretrained(clip_id).to("cuda")
text_encoder = CLIPTextModelWithProjection.from_pretrained(clip_id).to(device)
image_processor = CLIPImageProcessor.from_pretrained(clip_id)
image_encoder = CLIPVisionModelWithProjection.from_pretrained(clip_id).to("cuda")
image_encoder = CLIPVisionModelWithProjection.from_pretrained(clip_id).to(device)
```
Notice that we are using a particular CLIP checkpoint, i.e., `openai/clip-vit-large-patch14`. This is because the Stable Diffusion pre-training was performed with this CLIP variant. For more details, refer to the [documentation](https://huggingface.co/docs/transformers/model_doc/clip).
@@ -350,7 +350,7 @@ class DirectionalSimilarity(nn.Module):
def preprocess_image(self, image):
image = self.image_processor(image, return_tensors="pt")["pixel_values"]
return {"pixel_values": image.to("cuda")}
return {"pixel_values": image.to(device)}
def tokenize_text(self, text):
inputs = self.tokenizer(
@@ -360,7 +360,7 @@ class DirectionalSimilarity(nn.Module):
truncation=True,
return_tensors="pt",
)
return {"input_ids": inputs.input_ids.to("cuda")}
return {"input_ids": inputs.input_ids.to(device)}
def encode_image(self, image):
preprocessed_image = self.preprocess_image(image)
@@ -459,7 +459,6 @@ with ZipFile(local_filepath, "r") as zipper:
```python
from PIL import Image
import os
import numpy as np
dataset_path = "sample-imagenet-images"
image_paths = sorted([os.path.join(dataset_path, x) for x in os.listdir(dataset_path)])
@@ -478,7 +477,6 @@ Now that the images are loaded, let's apply some lightweight pre-processing on t
```python
from torchvision.transforms import functional as F
import torch
def preprocess_image(image):
@@ -500,10 +498,6 @@ dit_pipeline = DiTPipeline.from_pretrained("facebook/DiT-XL-2-256", torch_dtype=
dit_pipeline.scheduler = DPMSolverMultistepScheduler.from_config(dit_pipeline.scheduler.config)
dit_pipeline = dit_pipeline.to("cuda")
seed = 0
generator = torch.manual_seed(seed)
words = [
"cassette player",
"chainsaw",

View File

@@ -65,7 +65,7 @@ Pipelines are designed to be easy to use (therefore do not follow [*Simple over
The following design principles are followed:
- Pipelines follow the single-file policy. All pipelines can be found in individual directories under src/diffusers/pipelines. One pipeline folder corresponds to one diffusion paper/project/release. Multiple pipeline files can be gathered in one pipeline folder, as its done for [`src/diffusers/pipelines/stable-diffusion`](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines/stable_diffusion). If pipelines share similar functionality, one can make use of the [# Copied from mechanism](https://github.com/huggingface/diffusers/blob/125d783076e5bd9785beb05367a2d2566843a271/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_img2img.py#L251).
- Pipelines all inherit from [`DiffusionPipeline`].
- Every pipeline consists of different model and scheduler components, that are documented in the [`model_index.json` file](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5/blob/main/model_index.json), are accessible under the same name as attributes of the pipeline and can be shared between pipelines with [`DiffusionPipeline.components`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.components) function.
- Every pipeline consists of different model and scheduler components, that are documented in the [`model_index.json` file](https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/model_index.json), are accessible under the same name as attributes of the pipeline and can be shared between pipelines with [`DiffusionPipeline.components`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.components) function.
- Every pipeline should be loadable via the [`DiffusionPipeline.from_pretrained`](https://huggingface.co/docs/diffusers/main/en/api/diffusion_pipeline#diffusers.DiffusionPipeline.from_pretrained) function.
- Pipelines should be used **only** for inference.
- Pipelines should be very readable, self-explanatory, and easy to tweak.

View File

@@ -95,17 +95,17 @@ print(f"Model downloaded at {model_path}")
Once you have downloaded a snapshot of the model, you can test it using Apple's Python script.
```shell
python -m python_coreml_stable_diffusion.pipeline --prompt "a photo of an astronaut riding a horse on mars" -i ./models/coreml-stable-diffusion-v1-4_original_packages/original/packages -o </path/to/output/image> --compute-unit CPU_AND_GPU --seed 93
python -m python_coreml_stable_diffusion.pipeline --prompt "a photo of an astronaut riding a horse on mars" -i models/coreml-stable-diffusion-v1-4_original_packages -o </path/to/output/image> --compute-unit CPU_AND_GPU --seed 93
```
Pass the path of the downloaded checkpoint with `-i` flag to the script. `--compute-unit` indicates the hardware you want to allow for inference. It must be one of the following options: `ALL`, `CPU_AND_GPU`, `CPU_ONLY`, `CPU_AND_NE`. You may also provide an optional output path, and a seed for reproducibility.
The inference script assumes you're using the original version of the Stable Diffusion model, `CompVis/stable-diffusion-v1-4`. If you use another model, you *have* to specify its Hub id in the inference command line, using the `--model-version` option. This works for models already supported and custom models you trained or fine-tuned yourself.
For example, if you want to use [`stable-diffusion-v1-5/stable-diffusion-v1-5`](https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5):
For example, if you want to use [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5):
```shell
python -m python_coreml_stable_diffusion.pipeline --prompt "a photo of an astronaut riding a horse on mars" --compute-unit ALL -o output --seed 93 -i models/coreml-stable-diffusion-v1-5_original_packages --model-version stable-diffusion-v1-5/stable-diffusion-v1-5
python -m python_coreml_stable_diffusion.pipeline --prompt "a photo of an astronaut riding a horse on mars" --compute-unit ALL -o output --seed 93 -i models/coreml-stable-diffusion-v1-5_original_packages --model-version runwayml/stable-diffusion-v1-5
```
## Core ML inference in Swift

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