Compare commits

..

1 Commits

Author SHA1 Message Date
Dhruv Nair
327243fd95 update 2024-02-19 14:39:46 +00:00
619 changed files with 15865 additions and 65402 deletions

View File

@@ -66,32 +66,32 @@ body:
Questions on DiffusionPipeline (Saving, Loading, From pretrained, ...):
Questions on pipelines:
- Stable Diffusion @yiyixuxu @DN6 @sayakpaul
- Stable Diffusion XL @yiyixuxu @sayakpaul @DN6
- Kandinsky @yiyixuxu
- ControlNet @sayakpaul @yiyixuxu @DN6
- T2I Adapter @sayakpaul @yiyixuxu @DN6
- IF @DN6
- Text-to-Video / Video-to-Video @DN6 @sayakpaul
- Wuerstchen @DN6
- Stable Diffusion @yiyixuxu @DN6 @sayakpaul @patrickvonplaten
- Stable Diffusion XL @yiyixuxu @sayakpaul @DN6 @patrickvonplaten
- Kandinsky @yiyixuxu @patrickvonplaten
- ControlNet @sayakpaul @yiyixuxu @DN6 @patrickvonplaten
- T2I Adapter @sayakpaul @yiyixuxu @DN6 @patrickvonplaten
- IF @DN6 @patrickvonplaten
- Text-to-Video / Video-to-Video @DN6 @sayakpaul @patrickvonplaten
- Wuerstchen @DN6 @patrickvonplaten
- Other: @yiyixuxu @DN6
Questions on models:
- UNet @DN6 @yiyixuxu @sayakpaul
- VAE @sayakpaul @DN6 @yiyixuxu
- Transformers/Attention @DN6 @yiyixuxu @sayakpaul @DN6
- UNet @DN6 @yiyixuxu @sayakpaul @patrickvonplaten
- VAE @sayakpaul @DN6 @yiyixuxu @patrickvonplaten
- Transformers/Attention @DN6 @yiyixuxu @sayakpaul @DN6 @patrickvonplaten
Questions on Schedulers: @yiyixuxu
Questions on Schedulers: @yiyixuxu @patrickvonplaten
Questions on LoRA: @sayakpaul
Questions on LoRA: @sayakpaul @patrickvonplaten
Questions on Textual Inversion: @sayakpaul
Questions on Textual Inversion: @sayakpaul @patrickvonplaten
Questions on Training:
- DreamBooth @sayakpaul
- Text-to-Image Fine-tuning @sayakpaul
- Textual Inversion @sayakpaul
- ControlNet @sayakpaul
- DreamBooth @sayakpaul @patrickvonplaten
- Text-to-Image Fine-tuning @sayakpaul @patrickvonplaten
- Textual Inversion @sayakpaul @patrickvonplaten
- ControlNet @sayakpaul @patrickvonplaten
Questions on Tests: @DN6 @sayakpaul @yiyixuxu
@@ -99,7 +99,7 @@ body:
Questions on JAX- and MPS-related things: @pcuenca
Questions on audio pipelines: @DN6
Questions on audio pipelines: @DN6 @patrickvonplaten

View File

@@ -38,13 +38,13 @@ members/contributors who may be interested in your PR.
Core library:
- Schedulers: @yiyixuxu
- Pipelines: @sayakpaul @yiyixuxu @DN6
- Training examples: @sayakpaul
- Docs: @stevhliu and @sayakpaul
- Schedulers: @yiyixuxu and @patrickvonplaten
- Pipelines: @patrickvonplaten and @sayakpaul
- Training examples: @sayakpaul and @patrickvonplaten
- Docs: @stevhliu and @yiyixuxu
- JAX and MPS: @pcuenca
- Audio: @sanchit-gandhi
- General functionalities: @sayakpaul @yiyixuxu @DN6
- General functionalities: @patrickvonplaten and @sayakpaul
Integrations:

View File

@@ -1,7 +1,6 @@
name: Benchmarking tests
on:
workflow_dispatch:
schedule:
- cron: "30 1 1,15 * *" # every 2 weeks on the 1st and the 15th of every month at 1:30 AM
@@ -31,9 +30,9 @@ jobs:
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install pandas peft
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install pandas peft
- name: Environment
run: |
python utils/print_env.py

View File

@@ -1,58 +1,21 @@
name: Test, build, and push Docker images
name: Build Docker images (nightly)
on:
pull_request: # During PRs, we just check if the changes Dockerfiles can be successfully built
branches:
- main
paths:
- "docker/**"
workflow_dispatch:
schedule:
- cron: "0 0 * * *" # every day at midnight
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
group: docker-image-builds
cancel-in-progress: false
env:
REGISTRY: diffusers
CI_SLACK_CHANNEL: ${{ secrets.CI_DOCKER_CHANNEL }}
jobs:
test-build-docker-images:
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
if: github.event_name == 'pull_request'
steps:
- name: Set up Docker Buildx
uses: docker/setup-buildx-action@v1
- name: Check out code
uses: actions/checkout@v3
- name: Find Changed Dockerfiles
id: file_changes
uses: jitterbit/get-changed-files@v1
with:
format: 'space-delimited'
token: ${{ secrets.GITHUB_TOKEN }}
- name: Build Changed Docker Images
run: |
CHANGED_FILES="${{ steps.file_changes.outputs.all }}"
for FILE in $CHANGED_FILES; do
if [[ "$FILE" == docker/*Dockerfile ]]; then
DOCKER_PATH="${FILE%/Dockerfile}"
DOCKER_TAG=$(basename "$DOCKER_PATH")
echo "Building Docker image for $DOCKER_TAG"
docker build -t "$DOCKER_TAG" "$DOCKER_PATH"
fi
done
if: steps.file_changes.outputs.all != ''
build-docker-images:
runs-on: ubuntu-latest
build-and-push-docker-images:
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
if: github.event_name != 'pull_request'
permissions:
contents: read
packages: write
@@ -73,13 +36,13 @@ jobs:
steps:
- name: Checkout repository
uses: actions/checkout@v3
- name: Set up Docker Buildx
uses: docker/setup-buildx-action@v1
- name: Login to Docker Hub
uses: docker/login-action@v2
with:
username: ${{ env.REGISTRY }}
password: ${{ secrets.DOCKERHUB_TOKEN }}
- name: Build and push
uses: docker/build-push-action@v3
with:
@@ -87,27 +50,3 @@ jobs:
context: ./docker/${{ matrix.image-name }}
push: true
tags: ${{ env.REGISTRY }}/${{ matrix.image-name }}:latest
- name: Post to a Slack channel
id: slack
uses: slackapi/slack-github-action@6c661ce58804a1a20f6dc5fbee7f0381b469e001
with:
# Slack channel id, channel name, or user id to post message.
# See also: https://api.slack.com/methods/chat.postMessage#channels
channel-id: ${{ env.CI_SLACK_CHANNEL }}
# For posting a rich message using Block Kit
payload: |
{
"text": "${{ matrix.image-name }} Docker Image build result: ${{ job.status }}\n${{ github.event.head_commit.url }}",
"blocks": [
{
"type": "section",
"text": {
"type": "mrkdwn",
"text": "${{ matrix.image-name }} Docker Image build result: ${{ job.status }}\n${{ github.event.head_commit.url }}"
}
}
]
}
env:
SLACK_BOT_TOKEN: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}

View File

@@ -7,10 +7,6 @@ on:
- doc-builder*
- v*-release
- v*-patch
paths:
- "src/diffusers/**.py"
- "examples/**"
- "docs/**"
jobs:
build:

View File

@@ -2,10 +2,6 @@ name: Build PR Documentation
on:
pull_request:
paths:
- "src/diffusers/**.py"
- "examples/**"
- "docs/**"
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}

View File

@@ -1,7 +1,6 @@
name: Nightly and release tests on main/release branch
name: Nightly tests on main
on:
workflow_dispatch:
schedule:
- cron: "0 0 * * *" # every day at midnight
@@ -13,348 +12,106 @@ env:
PYTEST_TIMEOUT: 600
RUN_SLOW: yes
RUN_NIGHTLY: yes
PIPELINE_USAGE_CUTOFF: 5000
SLACK_API_TOKEN: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines Matrix
runs-on: ubuntu-latest
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
pip install -e .
pip install huggingface_hub
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
run: |
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
echo $matrix
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: test-pipelines.json
path: reports
run_nightly_tests_for_torch_pipelines:
name: Torch Pipelines CUDA Nightly Tests
needs: setup_torch_cuda_pipeline_matrix
run_nightly_tests:
strategy:
fail-fast: false
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on: [single-gpu, nvidia-gpu, t4, ci]
config:
- name: Nightly PyTorch CUDA tests on Ubuntu
framework: pytorch
runner: docker-gpu
image: diffusers/diffusers-pytorch-cuda
report: torch_cuda
- name: Nightly Flax TPU tests on Ubuntu
framework: flax
runner: docker-tpu
image: diffusers/diffusers-flax-tpu
report: flax_tpu
- name: Nightly ONNXRuntime CUDA tests on Ubuntu
framework: onnxruntime
runner: docker-gpu
image: diffusers/diffusers-onnxruntime-cuda
report: onnx_cuda
name: ${{ matrix.config.name }}
runs-on: ${{ matrix.config.runner }}
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
image: ${{ matrix.config.image }}
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ ${{ matrix.config.runner == 'docker-tpu' && '--privileged' || '--gpus 0'}}
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: nvidia-smi
if: ${{ matrix.config.runner == 'docker-gpu' }}
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
python -m pip install -e .[quality,test]
python -m pip install -U git+https://github.com/huggingface/transformers
python -m pip install git+https://github.com/huggingface/accelerate
- name: Environment
run: |
python utils/print_env.py
- name: Nightly PyTorch CUDA checkpoint (pipelines) tests
- name: Run nightly PyTorch CUDA tests
if: ${{ matrix.config.framework == 'pytorch' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
--report-log=tests_pipeline_${{ matrix.module }}_cuda.log \
tests/pipelines/${{ matrix.module }}
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Run nightly Flax TPU tests
if: ${{ matrix.config.framework == 'flax' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Run nightly ONNXRuntime CUDA tests
if: ${{ matrix.config.framework == 'onnxruntime' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
run: cat reports/tests_${{ matrix.config.report }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: pipeline_${{ matrix.module }}_test_reports
name: ${{ matrix.config.report }}_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_tests_for_other_torch_modules:
name: Torch Non-Pipelines CUDA Nightly Tests
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
strategy:
matrix:
module: [models, schedulers, others, examples]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly PyTorch CUDA tests for non-pipeline modules
if: ${{ matrix.module != 'examples'}}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_${{ matrix.module }}_cuda \
--report-log=tests_torch_${{ matrix.module }}_cuda.log \
tests/${{ matrix.module }}
- name: Run nightly example tests with Torch
if: ${{ matrix.module == 'examples' }}
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v --make-reports=examples_torch_cuda \
--report-log=examples_torch_cuda.log \
examples/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_torch_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_${{ matrix.module }}_cuda_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_lora_nightly_tests:
name: Nightly LoRA Tests with PEFT and TORCH
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly LoRA tests with PEFT and Torch
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_lora_cuda \
--report-log=tests_torch_lora_cuda.log \
tests/lora
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_lora_cuda_stats.txt
cat reports/tests_torch_lora_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: torch_lora_cuda_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_flax_tpu_tests:
name: Nightly Flax TPU Tests
runs-on: docker-tpu
if: github.event_name == 'schedule'
container:
image: diffusers/diffusers-flax-tpu
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --privileged
defaults:
run:
shell: bash
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly Flax TPU tests
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
--make-reports=tests_flax_tpu \
--report-log=tests_flax_tpu.log \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_flax_tpu_stats.txt
cat reports/tests_flax_tpu_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: flax_tpu_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_onnx_tests:
name: Nightly ONNXRuntime CUDA tests on Ubuntu
runs-on: docker-gpu
container:
image: diffusers/diffusers-onnxruntime-cuda
options: --gpus 0 --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m uv pip install pytest-reportlog
- name: Environment
run: python utils/print_env.py
- name: Run nightly ONNXRuntime CUDA tests
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_onnx_cuda \
--report-log=tests_onnx_cuda.log \
tests/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_onnx_cuda_stats.txt
cat reports/tests_onnx_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
with:
name: ${{ matrix.config.report }}_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_tests_apple_m1:
name: Nightly PyTorch MPS tests on MacOS
runs-on: [ self-hosted, apple-m1 ]
if: github.event_name == 'schedule'
steps:
- name: Checkout diffusers
@@ -375,11 +132,10 @@ jobs:
- name: Install dependencies
shell: arch -arch arm64 bash {0}
run: |
${CONDA_RUN} python -m pip install --upgrade pip uv
${CONDA_RUN} python -m uv pip install -e [quality,test]
${CONDA_RUN} python -m uv pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
${CONDA_RUN} python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate
${CONDA_RUN} python -m uv pip install pytest-reportlog
${CONDA_RUN} python -m pip install --upgrade pip
${CONDA_RUN} python -m pip install -e .[quality,test]
${CONDA_RUN} python -m pip install torch torchvision torchaudio --extra-index-url https://download.pytorch.org/whl/cpu
${CONDA_RUN} python -m pip install git+https://github.com/huggingface/accelerate
- name: Environment
shell: arch -arch arm64 bash {0}
@@ -392,9 +148,7 @@ jobs:
HF_HOME: /System/Volumes/Data/mnt/cache
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
--report-log=tests_torch_mps.log \
tests/
${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps tests/
- name: Failure short reports
if: ${{ failure() }}
@@ -406,9 +160,3 @@ jobs:
with:
name: torch_mps_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python scripts/log_reports.py >> $GITHUB_STEP_SUMMARY

View File

@@ -1,23 +0,0 @@
name: Notify Slack about a release
on:
workflow_dispatch:
release:
types: [published]
jobs:
build:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Setup Python
uses: actions/setup-python@v4
with:
python-version: '3.8'
- name: Notify Slack about the release
env:
SLACK_WEBHOOK_URL: ${{ secrets.SLACK_WEBHOOK_URL }}
run: pip install requests && python utils/notify_slack_about_release.py

View File

@@ -4,8 +4,6 @@ on:
pull_request:
branches:
- main
paths:
- "src/diffusers/**.py"
push:
branches:
- main
@@ -25,12 +23,10 @@ jobs:
python-version: "3.8"
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pip install --upgrade pip uv
python -m uv pip install -e .
python -m uv pip install pytest
python -m pip install --upgrade pip
pip install -e .
pip install pytest
- name: Check for soft dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
pytest tests/others/test_dependencies.py

View File

@@ -4,8 +4,6 @@ on:
pull_request:
branches:
- main
paths:
- "src/diffusers/**.py"
push:
branches:
- main
@@ -25,14 +23,12 @@ jobs:
python-version: "3.8"
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pip install --upgrade pip uv
python -m uv pip install -e .
python -m uv pip install "jax[cpu]>=0.2.16,!=0.3.2"
python -m uv pip install "flax>=0.4.1"
python -m uv pip install "jaxlib>=0.1.65"
python -m uv pip install pytest
python -m pip install --upgrade pip
pip install -e .
pip install "jax[cpu]>=0.2.16,!=0.3.2"
pip install "flax>=0.4.1"
pip install "jaxlib>=0.1.65"
pip install pytest
- name: Check for soft dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
pytest tests/others/test_dependencies.py

49
.github/workflows/pr_quality.yml vendored Normal file
View File

@@ -0,0 +1,49 @@
name: Run code quality checks
on:
pull_request:
branches:
- main
push:
branches:
- main
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
jobs:
check_code_quality:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check quality
run: |
ruff check examples tests src utils scripts
ruff format examples tests src utils scripts --check
check_repository_consistency:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check quality
run: |
python utils/check_copies.py
python utils/check_dummies.py
make deps_table_check_updated

View File

@@ -32,8 +32,8 @@ jobs:
fetch-depth: 0
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
- name: Environment
run: |
python utils/print_env.py
@@ -88,18 +88,16 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pip install -e [quality,test]
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install accelerate
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run all selected tests on CPU
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 2 --dist=loadfile -v --make-reports=${{ matrix.modules }}_tests_cpu ${{ fromJson(needs.setup_pr_tests.outputs.test_map)[matrix.modules] }}
- name: Failure short reports
@@ -145,18 +143,16 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pip install -e [quality,test]
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run Hub tests for models, schedulers, and pipelines on a staging env
if: ${{ matrix.config.framework == 'hub_tests_pytorch' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
HUGGINGFACE_CO_STAGING=true python -m pytest \
-m "is_staging_test" \
--make-reports=tests_${{ matrix.config.report }} \

View File

@@ -4,9 +4,6 @@ on:
pull_request:
branches:
- main
paths:
- "src/diffusers/**.py"
- "tests/**.py"
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
@@ -19,50 +16,7 @@ env:
PYTEST_TIMEOUT: 60
jobs:
check_code_quality:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check quality
run: make quality
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Quality check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make style && make quality'" >> $GITHUB_STEP_SUMMARY
check_repository_consistency:
needs: check_code_quality
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check repo consistency
run: |
python utils/check_copies.py
python utils/check_dummies.py
make deps_table_check_updated
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Repo consistency check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make fix-copies'" >> $GITHUB_STEP_SUMMARY
run_fast_tests:
needs: [check_code_quality, check_repository_consistency]
strategy:
fail-fast: false
matrix:
@@ -71,7 +25,7 @@ jobs:
name: LoRA - ${{ matrix.lib-versions }}
runs-on: [ self-hosted, intel-cpu, 8-cpu, ci ]
runs-on: docker-cpu
container:
image: diffusers/diffusers-pytorch-cpu
@@ -89,25 +43,23 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
if [ "${{ matrix.lib-versions }}" == "main" ]; then
python -m pip install -U peft@git+https://github.com/huggingface/peft.git
python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git
python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
python -m pip install -U git+https://github.com/huggingface/peft.git
python -m pip install -U git+https://github.com/huggingface/transformers.git
python -m pip install -U git+https://github.com/huggingface/accelerate.git
else
python -m uv pip install -U peft transformers accelerate
python -m pip install -U peft transformers accelerate
fi
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run fast PyTorch LoRA CPU tests with PEFT backend
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v \
--make-reports=tests_${{ matrix.config.report }} \
tests/lora/
tests/lora/test_lora_layers_peft.py

View File

@@ -4,14 +4,6 @@ on:
pull_request:
branches:
- main
paths:
- "src/diffusers/**.py"
- "benchmarks/**.py"
- "examples/**.py"
- "scripts/**.py"
- "tests/**.py"
- ".github/**.yml"
- "utils/**.py"
push:
branches:
- ci-*
@@ -27,72 +19,29 @@ env:
PYTEST_TIMEOUT: 60
jobs:
check_code_quality:
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check quality
run: make quality
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Quality check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make style && make quality'" >> $GITHUB_STEP_SUMMARY
check_repository_consistency:
needs: check_code_quality
runs-on: ubuntu-latest
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check repo consistency
run: |
python utils/check_copies.py
python utils/check_dummies.py
make deps_table_check_updated
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Repo consistency check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make fix-copies'" >> $GITHUB_STEP_SUMMARY
run_fast_tests:
needs: [check_code_quality, check_repository_consistency]
strategy:
fail-fast: false
matrix:
config:
- name: Fast PyTorch Pipeline CPU tests
framework: pytorch_pipelines
runner: [ self-hosted, intel-cpu, 32-cpu, 256-ram, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_pipelines
- name: Fast PyTorch Models & Schedulers CPU tests
framework: pytorch_models
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu_models_schedulers
- name: Fast Flax CPU tests
framework: flax
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-flax-cpu
report: flax_cpu
- name: PyTorch Example CPU tests
framework: pytorch_examples
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_example_cpu
@@ -116,20 +65,18 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install accelerate
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run fast PyTorch Pipeline CPU tests
if: ${{ matrix.config.framework == 'pytorch_pipelines' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 8 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/pipelines
@@ -137,8 +84,7 @@ jobs:
- name: Run fast PyTorch Model Scheduler CPU tests
if: ${{ matrix.config.framework == 'pytorch_models' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx and not Dependency" \
--make-reports=tests_${{ matrix.config.report }} \
tests/models tests/schedulers tests/others
@@ -146,8 +92,7 @@ jobs:
- name: Run fast Flax TPU tests
if: ${{ matrix.config.framework == 'flax' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Flax" \
--make-reports=tests_${{ matrix.config.report }} \
tests
@@ -155,9 +100,8 @@ jobs:
- name: Run example PyTorch CPU tests
if: ${{ matrix.config.framework == 'pytorch_examples' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install peft
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pip install peft
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_${{ matrix.config.report }} \
examples
@@ -173,14 +117,13 @@ jobs:
path: reports
run_staging_tests:
needs: [check_code_quality, check_repository_consistency]
strategy:
fail-fast: false
matrix:
config:
- name: Hub tests for models, schedulers, and pipelines
framework: hub_tests_pytorch
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_hub
@@ -204,18 +147,16 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run Hub tests for models, schedulers, and pipelines on a staging env
if: ${{ matrix.config.framework == 'hub_tests_pytorch' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
HUGGINGFACE_CO_STAGING=true python -m pytest \
-m "is_staging_test" \
--make-reports=tests_${{ matrix.config.report }} \

View File

@@ -4,8 +4,6 @@ on:
pull_request:
branches:
- main
paths:
- "src/diffusers/**.py"
push:
branches:
- main
@@ -25,12 +23,10 @@ jobs:
python-version: "3.8"
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pip install --upgrade pip uv
python -m uv pip install -e .
python -m uv pip install torch torchvision torchaudio
python -m uv pip install pytest
python -m pip install --upgrade pip
pip install -e .
pip install torch torchvision torchaudio
pip install pytest
- name: Check for soft dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
pytest tests/others/test_dependencies.py

View File

@@ -4,10 +4,7 @@ on:
push:
branches:
- main
paths:
- "src/diffusers/**.py"
- "examples/**.py"
- "tests/**.py"
env:
DIFFUSERS_IS_CI: yes
@@ -21,7 +18,10 @@ env:
jobs:
setup_torch_cuda_pipeline_matrix:
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on: ubuntu-latest
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cpu # this is a CPU image, but we need it to fetch the matrix
options: --shm-size "16gb" --ipc host
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
@@ -29,20 +29,23 @@ jobs:
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
pip install -e .
pip install huggingface_hub
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
run: |
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
echo $matrix
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v2
@@ -55,12 +58,13 @@ jobs:
needs: setup_torch_cuda_pipeline_matrix
strategy:
fail-fast: false
max-parallel: 1
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on: [single-gpu, nvidia-gpu, t4, ci]
runs-on: docker-gpu
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0 --privileged
options: --shm-size "16gb" --ipc host -v /mnt/hf_cache:/mnt/cache/ --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -69,17 +73,11 @@ jobs:
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Tailscale
uses: huggingface/tailscale-action@v1
with:
authkey: ${{ secrets.TAILSCALE_SSH_AUTHKEY }}
slackChannel: ${{ secrets.SLACK_CIFEEDBACK_CHANNEL }}
slackToken: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
python utils/print_env.py
@@ -93,12 +91,6 @@ jobs:
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
tests/pipelines/${{ matrix.module }}
- name: Tailscale Wait
if: ${{ failure() || runner.debug == '1' }}
uses: huggingface/tailscale-action@v1
with:
waitForSSH: true
authkey: ${{ secrets.TAILSCALE_SSH_AUTHKEY }}
- name: Failure short reports
if: ${{ failure() }}
run: |
@@ -132,9 +124,9 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
@@ -181,10 +173,10 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
python -m pip install -U peft@git+https://github.com/huggingface/peft.git
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
python -m pip install git+https://github.com/huggingface/peft.git
- name: Environment
run: |
@@ -231,9 +223,9 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
@@ -278,9 +270,9 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
python -m pip install git+https://github.com/huggingface/accelerate.git
- name: Environment
run: |
@@ -328,8 +320,7 @@ jobs:
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test,training]
python -m pip install -e .[quality,test,training]
- name: Environment
run: |
python utils/print_env.py
@@ -369,8 +360,7 @@ jobs:
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test,training]
python -m pip install -e .[quality,test,training]
- name: Environment
run: |
python utils/print_env.py
@@ -411,19 +401,16 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test,training]
python -m pip install -e .[quality,test,training]
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run example tests on GPU
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.HUGGING_FACE_HUB_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=examples_torch_cuda examples/
- name: Failure short reports
@@ -437,4 +424,4 @@ jobs:
uses: actions/upload-artifact@v2
with:
name: examples_test_reports
path: reports
path: reports

View File

@@ -4,10 +4,6 @@ on:
push:
branches:
- main
paths:
- "src/diffusers/**.py"
- "examples/**.py"
- "tests/**.py"
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
@@ -29,22 +25,22 @@ jobs:
config:
- name: Fast PyTorch CPU tests on Ubuntu
framework: pytorch
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_cpu
- name: Fast Flax CPU tests on Ubuntu
framework: flax
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-flax-cpu
report: flax_cpu
- name: Fast ONNXRuntime CPU tests on Ubuntu
framework: onnxruntime
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-onnxruntime-cpu
report: onnx_cpu
- name: PyTorch Example CPU tests on Ubuntu
framework: pytorch_examples
runner: [ self-hosted, intel-cpu, 8-cpu, ci ]
runner: docker-cpu
image: diffusers/diffusers-pytorch-cpu
report: torch_example_cpu
@@ -68,19 +64,17 @@ jobs:
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
apt-get update && apt-get install libsndfile1-dev libgl1 -y
python -m pip install -e .[quality,test]
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run fast PyTorch CPU tests
if: ${{ matrix.config.framework == 'pytorch' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
@@ -88,8 +82,7 @@ jobs:
- name: Run fast Flax TPU tests
if: ${{ matrix.config.framework == 'flax' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Flax" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
@@ -97,8 +90,7 @@ jobs:
- name: Run fast ONNXRuntime CPU tests
if: ${{ matrix.config.framework == 'onnxruntime' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
--make-reports=tests_${{ matrix.config.report }} \
tests/
@@ -106,9 +98,8 @@ jobs:
- name: Run example PyTorch CPU tests
if: ${{ matrix.config.framework == 'pytorch_examples' }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install peft
python -m pytest -n 4 --max-worker-restart=0 --dist=loadfile \
python -m pip install peft
python -m pytest -n 2 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_${{ matrix.config.report }} \
examples

View File

@@ -4,9 +4,6 @@ on:
push:
branches:
- main
paths:
- "src/diffusers/**.py"
- "tests/**.py"
env:
DIFFUSERS_IS_CI: yes
@@ -44,11 +41,11 @@ jobs:
- name: Install dependencies
shell: arch -arch arm64 bash {0}
run: |
${CONDA_RUN} python -m pip install --upgrade pip uv
${CONDA_RUN} python -m uv pip install -e [quality,test]
${CONDA_RUN} python -m uv pip install torch torchvision torchaudio
${CONDA_RUN} python -m uv pip install accelerate@git+https://github.com/huggingface/accelerate.git
${CONDA_RUN} python -m uv pip install transformers --upgrade
${CONDA_RUN} python -m pip install --upgrade pip
${CONDA_RUN} python -m pip install -e .[quality,test]
${CONDA_RUN} python -m pip install torch torchvision torchaudio
${CONDA_RUN} python -m pip install git+https://github.com/huggingface/accelerate.git
${CONDA_RUN} python -m pip install transformers --upgrade
- name: Environment
shell: arch -arch arm64 bash {0}

View File

@@ -1,81 +0,0 @@
# Adapted from https://blog.deepjyoti30.dev/pypi-release-github-action
name: PyPI release
on:
workflow_dispatch:
push:
tags:
- "*"
jobs:
find-and-checkout-latest-branch:
runs-on: ubuntu-latest
outputs:
latest_branch: ${{ steps.set_latest_branch.outputs.latest_branch }}
steps:
- name: Checkout Repo
uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: '3.8'
- name: Fetch latest branch
id: fetch_latest_branch
run: |
pip install -U requests packaging
LATEST_BRANCH=$(python utils/fetch_latest_release_branch.py)
echo "Latest branch: $LATEST_BRANCH"
echo "latest_branch=$LATEST_BRANCH" >> $GITHUB_ENV
- name: Set latest branch output
id: set_latest_branch
run: echo "::set-output name=latest_branch::${{ env.latest_branch }}"
release:
needs: find-and-checkout-latest-branch
runs-on: ubuntu-latest
steps:
- name: Checkout Repo
uses: actions/checkout@v3
with:
ref: ${{ needs.find-and-checkout-latest-branch.outputs.latest_branch }}
- name: Setup Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install -U setuptools wheel twine
pip install -U torch --index-url https://download.pytorch.org/whl/cpu
pip install -U transformers
- name: Build the dist files
run: python setup.py bdist_wheel && python setup.py sdist
- name: Publish to the test PyPI
env:
TWINE_USERNAME: ${{ secrets.TEST_PYPI_USERNAME }}
TWINE_PASSWORD: ${{ secrets.TEST_PYPI_PASSWORD }}
run: twine upload dist/* -r pypitest --repository-url=https://test.pypi.org/legacy/
- name: Test installing diffusers and importing
run: |
pip install diffusers && pip uninstall diffusers -y
pip install -i https://testpypi.python.org/pypi diffusers
python -c "from diffusers import __version__; print(__version__)"
python -c "from diffusers import DiffusionPipeline; pipe = DiffusionPipeline.from_pretrained('fusing/unet-ldm-dummy-update'); pipe()"
python -c "from diffusers import DiffusionPipeline; pipe = DiffusionPipeline.from_pretrained('hf-internal-testing/tiny-stable-diffusion-pipe', safety_checker=None); pipe('ah suh du')"
python -c "from diffusers import *"
- name: Publish to PyPI
env:
TWINE_USERNAME: ${{ secrets.PYPI_USERNAME }}
TWINE_PASSWORD: ${{ secrets.PYPI_PASSWORD }}
run: twine upload dist/* -r pypi

View File

@@ -1,46 +0,0 @@
name: SSH into runners
on:
workflow_dispatch:
inputs:
runner_type:
description: 'Type of runner to test (a10 or t4)'
required: true
docker_image:
description: 'Name of the Docker image'
required: true
env:
IS_GITHUB_CI: "1"
HF_HUB_READ_TOKEN: ${{ secrets.HF_HUB_READ_TOKEN }}
HF_HOME: /mnt/cache
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
RUN_SLOW: yes
jobs:
ssh_runner:
name: "SSH"
runs-on: [single-gpu, nvidia-gpu, "${{ github.event.inputs.runner_type }}", ci]
container:
image: ${{ github.event.inputs.docker_image }}
options: --gpus all --privileged --ipc host -v /mnt/cache/.cache/huggingface:/mnt/cache/
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Tailscale # In order to be able to SSH when a test fails
uses: huggingface/tailscale-action@v1
with:
authkey: ${{ secrets.TAILSCALE_SSH_AUTHKEY }}
slackChannel: ${{ secrets.SLACK_CIFEEDBACK_CHANNEL }}
slackToken: ${{ secrets.SLACK_CIFEEDBACK_BOT_TOKEN }}
waitForSSH: true

View File

@@ -1,30 +0,0 @@
name: Update Diffusers metadata
on:
workflow_dispatch:
push:
branches:
- main
- update_diffusers_metadata*
jobs:
update_metadata:
runs-on: ubuntu-22.04
defaults:
run:
shell: bash -l {0}
steps:
- uses: actions/checkout@v3
- name: Setup environment
run: |
pip install --upgrade pip
pip install datasets pandas
pip install .[torch]
- name: Update metadata
env:
HUGGING_FACE_HUB_TOKEN: ${{ secrets.SAYAK_HF_TOKEN }}
run: |
python utils/update_metadata.py --commit_sha ${{ github.sha }}

View File

@@ -19,16 +19,6 @@ authors:
family-names: Rasul
- given-names: Mishig
family-names: Davaadorj
- given-names: Dhruv
family-names: Nair
- given-names: Sayak
family-names: Paul
- given-names: Steven
family-names: Liu
- given-names: William
family-names: Berman
- given-names: Yiyi
family-names: Xu
- given-names: Thomas
family-names: Wolf
repository-code: 'https://github.com/huggingface/diffusers'

View File

@@ -42,7 +42,6 @@ repo-consistency:
quality:
ruff check $(check_dirs) setup.py
ruff format --check $(check_dirs) setup.py
doc-builder style src/diffusers docs/source --max_len 119 --check_only
python utils/check_doc_toc.py
# Format source code automatically and check is there are any problems left that need manual fixing
@@ -56,7 +55,6 @@ extra_style_checks:
style:
ruff check $(check_dirs) setup.py --fix
ruff format $(check_dirs) setup.py
doc-builder style src/diffusers docs/source --max_len 119
${MAKE} autogenerate_code
${MAKE} extra_style_checks

View File

@@ -77,7 +77,7 @@ Please refer to the [How to use Stable Diffusion in Apple Silicon](https://huggi
## Quickstart
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 22000+ checkpoints):
Generating outputs is super easy with 🤗 Diffusers. To generate an image from text, use the `from_pretrained` method to load any pretrained diffusion model (browse the [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) for 19000+ checkpoints):
```python
from diffusers import DiffusionPipeline
@@ -219,7 +219,7 @@ Also, say 👋 in our public Discord channel <a href="https://discord.gg/G7tWnz9
- https://github.com/deep-floyd/IF
- https://github.com/bentoml/BentoML
- https://github.com/bmaltais/kohya_ss
- +9000 other amazing GitHub repositories 💪
- +8000 other amazing GitHub repositories 💪
Thank you for using us ❤️.
@@ -238,7 +238,7 @@ We also want to thank @heejkoo for the very helpful overview of papers, code and
```bibtex
@misc{von-platen-etal-2022-diffusers,
author = {Patrick von Platen and Suraj Patil and Anton Lozhkov and Pedro Cuenca and Nathan Lambert and Kashif Rasul and Mishig Davaadorj and Dhruv Nair and Sayak Paul and William Berman and Yiyi Xu and Steven Liu and Thomas Wolf},
author = {Patrick von Platen and Suraj Patil and Anton Lozhkov and Pedro Cuenca and Nathan Lambert and Kashif Rasul and Mishig Davaadorj and Thomas Wolf},
title = {Diffusers: State-of-the-art diffusion models},
year = {2022},
publisher = {GitHub},

View File

@@ -141,7 +141,6 @@ class LCMLoRATextToImageBenchmark(TextToImageBenchmark):
super().__init__(args)
self.pipe.load_lora_weights(self.lora_id)
self.pipe.fuse_lora()
self.pipe.unload_lora_weights()
self.pipe.scheduler = LCMScheduler.from_config(self.pipe.scheduler.config)
def get_result_filepath(self, args):
@@ -236,35 +235,6 @@ class InpaintingBenchmark(ImageToImageBenchmark):
)
class IPAdapterTextToImageBenchmark(TextToImageBenchmark):
url = "https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/load_neg_embed.png"
image = load_image(url)
def __init__(self, args):
pipe = self.pipeline_class.from_pretrained(args.ckpt, torch_dtype=torch.float16).to("cuda")
pipe.load_ip_adapter(
args.ip_adapter_id[0],
subfolder="models" if "sdxl" not in args.ip_adapter_id[1] else "sdxl_models",
weight_name=args.ip_adapter_id[1],
)
if args.run_compile:
pipe.unet.to(memory_format=torch.channels_last)
print("Run torch compile")
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
pipe.set_progress_bar_config(disable=True)
self.pipe = pipe
def run_inference(self, pipe, args):
_ = pipe(
prompt=PROMPT,
ip_adapter_image=self.image,
num_inference_steps=args.num_inference_steps,
num_images_per_prompt=args.batch_size,
)
class ControlNetBenchmark(TextToImageBenchmark):
pipeline_class = StableDiffusionControlNetPipeline
aux_network_class = ControlNetModel

View File

@@ -1,32 +0,0 @@
import argparse
import sys
sys.path.append(".")
from base_classes import IPAdapterTextToImageBenchmark # noqa: E402
IP_ADAPTER_CKPTS = {
"runwayml/stable-diffusion-v1-5": ("h94/IP-Adapter", "ip-adapter_sd15.bin"),
"stabilityai/stable-diffusion-xl-base-1.0": ("h94/IP-Adapter", "ip-adapter_sdxl.bin"),
}
if __name__ == "__main__":
parser = argparse.ArgumentParser()
parser.add_argument(
"--ckpt",
type=str,
default="runwayml/stable-diffusion-v1-5",
choices=list(IP_ADAPTER_CKPTS.keys()),
)
parser.add_argument("--batch_size", type=int, default=1)
parser.add_argument("--num_inference_steps", type=int, default=50)
parser.add_argument("--model_cpu_offload", action="store_true")
parser.add_argument("--run_compile", action="store_true")
args = parser.parse_args()
args.ip_adapter_id = IP_ADAPTER_CKPTS[args.ckpt]
benchmark_pipe = IPAdapterTextToImageBenchmark(args)
args.ckpt = f"{args.ckpt} (IP-Adapter)"
benchmark_pipe.benchmark(args)

View File

@@ -72,7 +72,7 @@ def main():
command += " --run_compile"
run_command(command.split())
elif file in ["benchmark_sd_inpainting.py", "benchmark_ip_adapters.py"]:
elif file == "benchmark_sd_inpainting.py":
sdxl_ckpt = "stabilityai/stable-diffusion-xl-base-1.0"
command = f"python {file} --ckpt {sdxl_ckpt}"
run_command(command.split())

View File

@@ -12,7 +12,6 @@ RUN apt update && \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -24,13 +23,13 @@ ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
# follow the instructions here: https://cloud.google.com/tpu/docs/run-in-container#train_a_jax_model_in_a_docker_container
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m uv pip install --upgrade --no-cache-dir \
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
python3 -m pip install --upgrade --no-cache-dir \
clu \
"jax[cpu]>=0.2.16,!=0.3.2" \
"flax>=0.4.1" \
"jaxlib>=0.1.65" && \
python3 -m uv pip install --no-cache-dir \
python3 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \

View File

@@ -12,7 +12,6 @@ RUN apt update && \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -24,15 +23,15 @@ ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
# follow the instructions here: https://cloud.google.com/tpu/docs/run-in-container#train_a_jax_model_in_a_docker_container
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
python3 -m pip install --no-cache-dir \
"jax[tpu]>=0.2.16,!=0.3.2" \
-f https://storage.googleapis.com/jax-releases/libtpu_releases.html && \
python3 -m uv pip install --upgrade --no-cache-dir \
python3 -m pip install --upgrade --no-cache-dir \
clu \
"flax>=0.4.1" \
"jaxlib>=0.1.65" && \
python3 -m uv pip install --no-cache-dir \
python3 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \

View File

@@ -12,7 +12,6 @@ RUN apt update && \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -23,14 +22,14 @@ RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m uv pip install --no-cache-dir \
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
python3 -m pip install --no-cache-dir \
torch==2.1.2 \
torchvision==0.16.2 \
torchaudio==2.1.2 \
onnxruntime \
--extra-index-url https://download.pytorch.org/whl/cpu && \
python3 -m uv pip install --no-cache-dir \
python3 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \

View File

@@ -1,4 +1,4 @@
FROM nvidia/cuda:12.1.0-runtime-ubuntu20.04
FROM nvidia/cuda:11.6.2-cudnn8-devel-ubuntu20.04
LABEL maintainer="Hugging Face"
LABEL repository="diffusers"
@@ -12,7 +12,6 @@ RUN apt update && \
curl \
ca-certificates \
libsndfile1-dev \
libgl1 \
python3.8 \
python3-pip \
python3.8-venv && \
@@ -23,14 +22,14 @@ RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m uv pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
python3 -m pip install --no-cache-dir \
torch==2.1.2 \
torchvision==0.16.2 \
torchaudio==2.1.2 \
"onnxruntime-gpu>=1.13.1" \
--extra-index-url https://download.pytorch.org/whl/cu117 && \
python3 -m uv pip install --no-cache-dir \
python3 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \

View File

@@ -24,8 +24,8 @@ RUN python3.9 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3.9 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3.9 -m uv pip install --no-cache-dir \
RUN python3.9 -m pip install --no-cache-dir --upgrade pip && \
python3.9 -m pip install --no-cache-dir \
torch \
torchvision \
torchaudio \

View File

@@ -23,14 +23,14 @@ RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m uv pip install --no-cache-dir \
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
python3 -m pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark \
--extra-index-url https://download.pytorch.org/whl/cpu && \
python3 -m uv pip install --no-cache-dir \
python3 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \
@@ -40,6 +40,6 @@ RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
numpy \
scipy \
tensorboard \
transformers matplotlib
transformers
CMD ["/bin/bash"]

View File

@@ -23,8 +23,8 @@ RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m uv pip install --no-cache-dir \
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
python3 -m pip install --no-cache-dir \
torch \
torchvision \
torchaudio \

View File

@@ -23,13 +23,13 @@ RUN python3 -m venv /opt/venv
ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
RUN python3 -m pip install --no-cache-dir --upgrade pip && \
python3 -m pip install --no-cache-dir \
torch \
torchvision \
torchaudio \
invisible_watermark && \
python3 -m uv pip install --no-cache-dir \
python3 -m pip install --no-cache-dir \
accelerate \
datasets \
hf-doc-builder \

View File

@@ -18,143 +18,155 @@
- local: tutorials/basic_training
title: Train a diffusion model
- local: tutorials/using_peft_for_inference
title: Load LoRAs for inference
title: Inference with PEFT
- local: tutorials/fast_diffusion
title: Accelerate inference of text-to-image diffusion models
title: Tutorials
- sections:
- local: using-diffusers/loading
title: Load pipelines
- local: using-diffusers/custom_pipeline_overview
title: Load community pipelines and components
- local: using-diffusers/schedulers
title: Load schedulers and models
- local: using-diffusers/using_safetensors
title: Load safetensors
- local: using-diffusers/other-formats
title: Load different Stable Diffusion formats
- local: using-diffusers/loading_adapters
title: Load adapters
- local: using-diffusers/push_to_hub
title: Push files to the Hub
title: Load pipelines and adapters
- sections:
- local: using-diffusers/unconditional_image_generation
title: Unconditional image generation
- local: using-diffusers/conditional_image_generation
title: Text-to-image
- local: using-diffusers/img2img
title: Image-to-image
- local: using-diffusers/inpaint
title: Inpainting
- local: using-diffusers/text-img2vid
title: Text or image-to-video
- local: using-diffusers/depth2img
title: Depth-to-image
title: Generative tasks
- sections:
- local: using-diffusers/overview_techniques
title: Overview
- local: training/distributed_inference
title: Distributed inference with multiple GPUs
- local: using-diffusers/merge_loras
title: Merge LoRAs
- local: using-diffusers/callback
title: Pipeline callbacks
- local: using-diffusers/reusing_seeds
title: Reproducible pipelines
- local: using-diffusers/image_quality
title: Controlling image quality
- local: using-diffusers/weighted_prompts
title: Prompt techniques
title: Inference techniques
- sections:
- local: using-diffusers/sdxl
title: Stable Diffusion XL
- local: using-diffusers/sdxl_turbo
title: SDXL Turbo
- local: using-diffusers/kandinsky
title: Kandinsky
- local: using-diffusers/ip_adapter
title: IP-Adapter
- local: using-diffusers/controlnet
title: ControlNet
- local: using-diffusers/t2i_adapter
title: T2I-Adapter
- local: using-diffusers/textual_inversion_inference
title: Textual inversion
- local: using-diffusers/shap-e
title: Shap-E
- local: using-diffusers/diffedit
title: DiffEdit
- local: using-diffusers/inference_with_lcm_lora
title: Latent Consistency Model-LoRA
- local: using-diffusers/inference_with_lcm
title: Latent Consistency Model
- local: using-diffusers/inference_with_tcd_lora
title: Trajectory Consistency Distillation-LoRA
- local: using-diffusers/svd
title: Stable Video Diffusion
title: Specific pipeline examples
- sections:
- local: training/overview
title: Overview
- local: training/create_dataset
title: Create a dataset for training
- local: training/adapt_a_model
title: Adapt a model to a new task
- sections:
- local: training/unconditional_training
- local: using-diffusers/loading_overview
title: Overview
- local: using-diffusers/loading
title: Load pipelines, models, and schedulers
- local: using-diffusers/schedulers
title: Load and compare different schedulers
- local: using-diffusers/custom_pipeline_overview
title: Load community pipelines and components
- local: using-diffusers/using_safetensors
title: Load safetensors
- local: using-diffusers/other-formats
title: Load different Stable Diffusion formats
- local: using-diffusers/loading_adapters
title: Load adapters
- local: using-diffusers/push_to_hub
title: Push files to the Hub
title: Loading & Hub
- sections:
- local: using-diffusers/pipeline_overview
title: Overview
- local: using-diffusers/unconditional_image_generation
title: Unconditional image generation
- local: training/text2image
- local: using-diffusers/conditional_image_generation
title: Text-to-image
- local: training/sdxl
title: Stable Diffusion XL
- local: training/kandinsky
title: Kandinsky 2.2
- local: training/wuerstchen
title: Wuerstchen
- local: training/controlnet
title: ControlNet
- local: training/t2i_adapters
title: T2I-Adapters
- local: training/instructpix2pix
title: InstructPix2Pix
title: Models
isExpanded: false
- local: using-diffusers/img2img
title: Image-to-image
- local: using-diffusers/inpaint
title: Inpainting
- local: using-diffusers/text-img2vid
title: Text or image-to-video
- local: using-diffusers/depth2img
title: Depth-to-image
title: Tasks
- sections:
- local: training/text_inversion
title: Textual Inversion
- local: training/dreambooth
title: DreamBooth
- local: training/lora
title: LoRA
- local: training/custom_diffusion
title: Custom Diffusion
- local: training/lcm_distill
title: Latent Consistency Distillation
- local: training/ddpo
title: Reinforcement learning training with DDPO
title: Methods
isExpanded: false
title: Training
- local: using-diffusers/textual_inversion_inference
title: Textual inversion
- local: using-diffusers/ip_adapter
title: IP-Adapter
- local: training/distributed_inference
title: Distributed inference with multiple GPUs
- local: using-diffusers/reusing_seeds
title: Improve image quality with deterministic generation
- local: using-diffusers/control_brightness
title: Control image brightness
- local: using-diffusers/weighted_prompts
title: Prompt weighting
- local: using-diffusers/freeu
title: Improve generation quality with FreeU
title: Techniques
- sections:
- local: using-diffusers/pipeline_overview
title: Overview
- local: using-diffusers/sdxl
title: Stable Diffusion XL
- local: using-diffusers/sdxl_turbo
title: SDXL Turbo
- local: using-diffusers/kandinsky
title: Kandinsky
- local: using-diffusers/controlnet
title: ControlNet
- local: using-diffusers/shap-e
title: Shap-E
- local: using-diffusers/diffedit
title: DiffEdit
- local: using-diffusers/distilled_sd
title: Distilled Stable Diffusion inference
- local: using-diffusers/callback
title: Pipeline callbacks
- local: using-diffusers/reproducibility
title: Create reproducible pipelines
- local: using-diffusers/custom_pipeline_examples
title: Community pipelines
- local: using-diffusers/contribute_pipeline
title: Contribute a community pipeline
- local: using-diffusers/inference_with_lcm_lora
title: Latent Consistency Model-LoRA
- local: using-diffusers/inference_with_lcm
title: Latent Consistency Model
- local: using-diffusers/svd
title: Stable Video Diffusion
title: Specific pipeline examples
- sections:
- local: training/overview
title: Overview
- local: training/create_dataset
title: Create a dataset for training
- local: training/adapt_a_model
title: Adapt a model to a new task
- sections:
- local: training/unconditional_training
title: Unconditional image generation
- local: training/text2image
title: Text-to-image
- local: training/sdxl
title: Stable Diffusion XL
- local: training/kandinsky
title: Kandinsky 2.2
- local: training/wuerstchen
title: Wuerstchen
- local: training/controlnet
title: ControlNet
- local: training/t2i_adapters
title: T2I-Adapters
- local: training/instructpix2pix
title: InstructPix2Pix
title: Models
- sections:
- local: training/text_inversion
title: Textual Inversion
- local: training/dreambooth
title: DreamBooth
- local: training/lora
title: LoRA
- local: training/custom_diffusion
title: Custom Diffusion
- local: training/lcm_distill
title: Latent Consistency Distillation
- local: training/ddpo
title: Reinforcement learning training with DDPO
title: Methods
title: Training
- sections:
- local: using-diffusers/other-modalities
title: Other Modalities
title: Taking Diffusers Beyond Images
title: Using Diffusers
- sections:
- local: optimization/fp16
title: Speed up inference
- local: using-diffusers/distilled_sd
title: Distilled Stable Diffusion inference
- local: optimization/memory
title: Reduce memory usage
- local: optimization/torch2.0
title: PyTorch 2.0
- local: optimization/xformers
title: xFormers
- local: optimization/tome
title: Token merging
- local: optimization/deepcache
title: DeepCache
- local: optimization/tgate
title: TGATE
- local: optimization/opt_overview
title: Overview
- sections:
- local: optimization/fp16
title: Speed up inference
- local: optimization/memory
title: Reduce memory usage
- local: optimization/torch2.0
title: PyTorch 2.0
- local: optimization/xformers
title: xFormers
- local: optimization/tome
title: Token merging
- local: optimization/deepcache
title: DeepCache
title: General optimizations
- sections:
- local: using-diffusers/stable_diffusion_jax_how_to
title: JAX/Flax
@@ -164,14 +176,14 @@
title: OpenVINO
- local: optimization/coreml
title: Core ML
title: Optimized model formats
title: Optimized model types
- sections:
- local: optimization/mps
title: Metal Performance Shaders (MPS)
- local: optimization/habana
title: Habana Gaudi
title: Optimized hardware
title: Accelerate inference and reduce memory
title: Optimization
- sections:
- local: conceptual/philosophy
title: Philosophy
@@ -193,7 +205,6 @@
- local: api/outputs
title: Outputs
title: Main Classes
isExpanded: false
- sections:
- local: api/loaders/ip_adapter
title: IP-Adapter
@@ -208,7 +219,6 @@
- local: api/loaders/peft
title: PEFT
title: Loaders
isExpanded: false
- sections:
- local: api/models/overview
title: Overview
@@ -243,7 +253,6 @@
- local: api/models/controlnet
title: ControlNet
title: Models
isExpanded: false
- sections:
- local: api/pipelines/overview
title: Overview
@@ -267,10 +276,6 @@
title: ControlNet
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/controlnetxs
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
title: ControlNet-XS with Stable Diffusion XL
- local: api/pipelines/dance_diffusion
title: Dance Diffusion
- local: api/pipelines/ddim
@@ -297,8 +302,6 @@
title: Latent Consistency Models
- local: api/pipelines/latent_diffusion
title: Latent Diffusion
- local: api/pipelines/ledits_pp
title: LEDITS++
- local: api/pipelines/panorama
title: MultiDiffusion
- local: api/pipelines/musicldm
@@ -315,8 +318,6 @@
title: Semantic Guidance
- local: api/pipelines/shap_e
title: Shap-E
- local: api/pipelines/stable_cascade
title: Stable Cascade
- sections:
- local: api/pipelines/stable_diffusion/overview
title: Overview
@@ -349,7 +350,7 @@
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D Upscaler
- local: api/pipelines/stable_diffusion/adapter
title: T2I-Adapter
title: Stable Diffusion T2I-Adapter
- local: api/pipelines/stable_diffusion/gligen
title: GLIGEN (Grounded Language-to-Image Generation)
title: Stable Diffusion
@@ -368,7 +369,6 @@
- local: api/pipelines/wuerstchen
title: Wuerstchen
title: Pipelines
isExpanded: false
- sections:
- local: api/schedulers/overview
title: Overview
@@ -392,10 +392,6 @@
title: DPMSolverSDEScheduler
- local: api/schedulers/singlestep_dpm_solver
title: DPMSolverSinglestepScheduler
- local: api/schedulers/edm_multistep_dpm_solver
title: EDMDPMSolverMultistepScheduler
- local: api/schedulers/edm_euler
title: EDMEulerScheduler
- local: api/schedulers/euler_ancestral
title: EulerAncestralDiscreteScheduler
- local: api/schedulers/euler
@@ -422,14 +418,11 @@
title: ScoreSdeVeScheduler
- local: api/schedulers/score_sde_vp
title: ScoreSdeVpScheduler
- local: api/schedulers/tcd
title: TCDScheduler
- local: api/schedulers/unipc
title: UniPCMultistepScheduler
- local: api/schedulers/vq_diffusion
title: VQDiffusionScheduler
title: Schedulers
isExpanded: false
- sections:
- local: api/internal_classes_overview
title: Overview
@@ -444,5 +437,4 @@
- local: api/image_processor
title: VAE Image Processor
title: Internal classes
isExpanded: false
title: API

View File

@@ -41,6 +41,12 @@ An attention processor is a class for applying different types of attention mech
## FusedAttnProcessor2_0
[[autodoc]] models.attention_processor.FusedAttnProcessor2_0
## LoRAAttnProcessor
[[autodoc]] models.attention_processor.LoRAAttnProcessor
## LoRAAttnProcessor2_0
[[autodoc]] models.attention_processor.LoRAAttnProcessor2_0
## LoRAAttnAddedKVProcessor
[[autodoc]] models.attention_processor.LoRAAttnAddedKVProcessor

View File

@@ -23,7 +23,3 @@ Learn how to load an IP-Adapter checkpoint and image in the IP-Adapter [loading]
## IPAdapterMixin
[[autodoc]] loaders.ip_adapter.IPAdapterMixin
## IPAdapterMaskProcessor
[[autodoc]] image_processor.IPAdapterMaskProcessor

View File

@@ -1,18 +1,6 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Consistency Decoder
Consistency decoder can be used to decode the latents from the denoising UNet in the [`StableDiffusionPipeline`]. This decoder was introduced in the [DALL-E 3 technical report](https://openai.com/dall-e-3).
Consistency decoder can be used to decode the latents from the denoising UNet in the [`StableDiffusionPipeline`]. This decoder was introduced in the [DALL-E 3 technical report](https://openai.com/dall-e-3).
The original codebase can be found at [openai/consistencydecoder](https://github.com/openai/consistencydecoder).

View File

@@ -408,114 +408,6 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
</Tip>
<table>
<tr>
<th align=center>Without FreeInit enabled</th>
<th align=center>With FreeInit enabled</th>
</tr>
<tr>
<td align=center>
panda playing a guitar
<br />
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-no-freeinit.gif"
alt="panda playing a guitar"
style="width: 300px;" />
</td>
<td align=center>
panda playing a guitar
<br/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatediff-freeinit.gif"
alt="panda playing a guitar"
style="width: 300px;" />
</td>
</tr>
</table>
## Using AnimateLCM
[AnimateLCM](https://animatelcm.github.io/) is a motion module checkpoint and an [LCM LoRA](https://huggingface.co/docs/diffusers/using-diffusers/inference_with_lcm_lora) that have been created using a consistency learning strategy that decouples the distillation of the image generation priors and the motion generation priors.
```python
import torch
from diffusers import AnimateDiffPipeline, LCMScheduler, MotionAdapter
from diffusers.utils import export_to_gif
adapter = MotionAdapter.from_pretrained("wangfuyun/AnimateLCM")
pipe = AnimateDiffPipeline.from_pretrained("emilianJR/epiCRealism", motion_adapter=adapter)
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config, beta_schedule="linear")
pipe.load_lora_weights("wangfuyun/AnimateLCM", weight_name="sd15_lora_beta.safetensors", adapter_name="lcm-lora")
pipe.enable_vae_slicing()
pipe.enable_model_cpu_offload()
output = pipe(
prompt="A space rocket with trails of smoke behind it launching into space from the desert, 4k, high resolution",
negative_prompt="bad quality, worse quality, low resolution",
num_frames=16,
guidance_scale=1.5,
num_inference_steps=6,
generator=torch.Generator("cpu").manual_seed(0),
)
frames = output.frames[0]
export_to_gif(frames, "animatelcm.gif")
```
<table>
<tr>
<td><center>
A space rocket, 4K.
<br>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatelcm-output.gif"
alt="A space rocket, 4K"
style="width: 300px;" />
</center></td>
</tr>
</table>
AnimateLCM is also compatible with existing [Motion LoRAs](https://huggingface.co/collections/dn6/animatediff-motion-loras-654cb8ad732b9e3cf4d3c17e).
```python
import torch
from diffusers import AnimateDiffPipeline, LCMScheduler, MotionAdapter
from diffusers.utils import export_to_gif
adapter = MotionAdapter.from_pretrained("wangfuyun/AnimateLCM")
pipe = AnimateDiffPipeline.from_pretrained("emilianJR/epiCRealism", motion_adapter=adapter)
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config, beta_schedule="linear")
pipe.load_lora_weights("wangfuyun/AnimateLCM", weight_name="sd15_lora_beta.safetensors", adapter_name="lcm-lora")
pipe.load_lora_weights("guoyww/animatediff-motion-lora-tilt-up", adapter_name="tilt-up")
pipe.set_adapters(["lcm-lora", "tilt-up"], [1.0, 0.8])
pipe.enable_vae_slicing()
pipe.enable_model_cpu_offload()
output = pipe(
prompt="A space rocket with trails of smoke behind it launching into space from the desert, 4k, high resolution",
negative_prompt="bad quality, worse quality, low resolution",
num_frames=16,
guidance_scale=1.5,
num_inference_steps=6,
generator=torch.Generator("cpu").manual_seed(0),
)
frames = output.frames[0]
export_to_gif(frames, "animatelcm-motion-lora.gif")
```
<table>
<tr>
<td><center>
A space rocket, 4K.
<br>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/animatelcm-motion-lora.gif"
alt="A space rocket, 4K"
style="width: 300px;" />
</center></td>
</tr>
</table>
## AnimateDiffPipeline
[[autodoc]] AnimateDiffPipeline

View File

@@ -20,8 +20,7 @@ The abstract of the paper is the following:
*Although audio generation shares commonalities across different types of audio, such as speech, music, and sound effects, designing models for each type requires careful consideration of specific objectives and biases that can significantly differ from those of other types. To bring us closer to a unified perspective of audio generation, this paper proposes a framework that utilizes the same learning method for speech, music, and sound effect generation. Our framework introduces a general representation of audio, called "language of audio" (LOA). Any audio can be translated into LOA based on AudioMAE, a self-supervised pre-trained representation learning model. In the generation process, we translate any modalities into LOA by using a GPT-2 model, and we perform self-supervised audio generation learning with a latent diffusion model conditioned on LOA. The proposed framework naturally brings advantages such as in-context learning abilities and reusable self-supervised pretrained AudioMAE and latent diffusion models. Experiments on the major benchmarks of text-to-audio, text-to-music, and text-to-speech demonstrate state-of-the-art or competitive performance against previous approaches. Our code, pretrained model, and demo are available at [this https URL](https://audioldm.github.io/audioldm2).*
This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi) and [Nguyễn Công Tú Anh](https://github.com/tuanh123789). The original codebase can be
found at [haoheliu/audioldm2](https://github.com/haoheliu/audioldm2).
This pipeline was contributed by [sanchit-gandhi](https://huggingface.co/sanchit-gandhi). The original codebase can be found at [haoheliu/audioldm2](https://github.com/haoheliu/audioldm2).
## Tips
@@ -37,8 +36,6 @@ See table below for details on the three checkpoints:
| [audioldm2](https://huggingface.co/cvssp/audioldm2) | Text-to-audio | 350M | 1.1B | 1150k |
| [audioldm2-large](https://huggingface.co/cvssp/audioldm2-large) | Text-to-audio | 750M | 1.5B | 1150k |
| [audioldm2-music](https://huggingface.co/cvssp/audioldm2-music) | Text-to-music | 350M | 1.1B | 665k |
| [audioldm2-gigaspeech](https://huggingface.co/anhnct/audioldm2_gigaspeech) | Text-to-speech | 350M | 1.1B |10k |
| [audioldm2-ljspeech](https://huggingface.co/anhnct/audioldm2_ljspeech) | Text-to-speech | 350M | 1.1B | |
### Constructing a prompt
@@ -56,7 +53,7 @@ See table below for details on the three checkpoints:
* The quality of the generated waveforms can vary significantly based on the seed. Try generating with different seeds until you find a satisfactory generation.
* Multiple waveforms can be generated in one go: set `num_waveforms_per_prompt` to a value greater than 1. Automatic scoring will be performed between the generated waveforms and prompt text, and the audios ranked from best to worst accordingly.
The following example demonstrates how to construct good music and speech generation using the aforementioned tips: [example](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2Pipeline.__call__.example).
The following example demonstrates how to construct good music generation using the aforementioned tips: [example](https://huggingface.co/docs/diffusers/main/en/api/pipelines/audioldm2#diffusers.AudioLDM2Pipeline.__call__.example).
<Tip>

View File

@@ -12,10 +12,42 @@ specific language governing permissions and limitations under the License.
# AutoPipeline
The `AutoPipeline` is designed to make it easy to load a checkpoint for a task without needing to know the specific pipeline class. Based on the task, the `AutoPipeline` automatically retrieves the correct pipeline class from the checkpoint `model_index.json` file.
`AutoPipeline` is designed to:
1. make it easy for you to load a checkpoint for a task without knowing the specific pipeline class to use
2. use multiple pipelines in your workflow
Based on the task, the `AutoPipeline` class automatically retrieves the relevant pipeline given the name or path to the pretrained weights with the `from_pretrained()` method.
To seamlessly switch between tasks with the same checkpoint without reallocating additional memory, use the `from_pipe()` method to transfer the components from the original pipeline to the new one.
```py
from diffusers import AutoPipelineForText2Image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
image = pipeline(prompt, num_inference_steps=25).images[0]
```
<Tip>
Check out the [AutoPipeline](../../tutorials/autopipeline) tutorial to learn how to use this API!
</Tip>
`AutoPipeline` supports text-to-image, image-to-image, and inpainting for the following diffusion models:
- [Stable Diffusion](./stable_diffusion/overview)
- [ControlNet](./controlnet)
- [Stable Diffusion XL (SDXL)](./stable_diffusion/stable_diffusion_xl)
- [DeepFloyd IF](./deepfloyd_if)
- [Kandinsky 2.1](./kandinsky)
- [Kandinsky 2.2](./kandinsky_v22)
> [!TIP]
> Check out the [AutoPipeline](../../tutorials/autopipeline) tutorial to learn how to use this API!
## AutoPipelineForText2Image

View File

@@ -1,54 +0,0 @@
<!--Copyright 2023 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# LEDITS++
LEDITS++ was proposed in [LEDITS++: Limitless Image Editing using Text-to-Image Models](https://huggingface.co/papers/2311.16711) by Manuel Brack, Felix Friedrich, Katharina Kornmeier, Linoy Tsaban, Patrick Schramowski, Kristian Kersting, Apolinário Passos.
The abstract from the paper is:
*Text-to-image diffusion models have recently received increasing interest for their astonishing ability to produce high-fidelity images from solely text inputs. Subsequent research efforts aim to exploit and apply their capabilities to real image editing. However, existing image-to-image methods are often inefficient, imprecise, and of limited versatility. They either require time-consuming fine-tuning, deviate unnecessarily strongly from the input image, and/or lack support for multiple, simultaneous edits. To address these issues, we introduce LEDITS++, an efficient yet versatile and precise textual image manipulation technique. LEDITS++'s novel inversion approach requires no tuning nor optimization and produces high-fidelity results with a few diffusion steps. Second, our methodology supports multiple simultaneous edits and is architecture-agnostic. Third, we use a novel implicit masking technique that limits changes to relevant image regions. We propose the novel TEdBench++ benchmark as part of our exhaustive evaluation. Our results demonstrate the capabilities of LEDITS++ and its improvements over previous methods. The project page is available at https://leditsplusplus-project.static.hf.space .*
<Tip>
You can find additional information about LEDITS++ on the [project page](https://leditsplusplus-project.static.hf.space/index.html) and try it out in a [demo](https://huggingface.co/spaces/editing-images/leditsplusplus).
</Tip>
<Tip warning={true}>
Due to some backward compatability issues with the current diffusers implementation of [`~schedulers.DPMSolverMultistepScheduler`] this implementation of LEdits++ can no longer guarantee perfect inversion.
This issue is unlikely to have any noticeable effects on applied use-cases. However, we provide an alternative implementation that guarantees perfect inversion in a dedicated [GitHub repo](https://github.com/ml-research/ledits_pp).
</Tip>
We provide two distinct pipelines based on different pre-trained models.
## LEditsPPPipelineStableDiffusion
[[autodoc]] pipelines.ledits_pp.LEditsPPPipelineStableDiffusion
- all
- __call__
- invert
## LEditsPPPipelineStableDiffusionXL
[[autodoc]] pipelines.ledits_pp.LEditsPPPipelineStableDiffusionXL
- all
- __call__
- invert
## LEditsPPDiffusionPipelineOutput
[[autodoc]] pipelines.ledits_pp.pipeline_output.LEditsPPDiffusionPipelineOutput
- all
## LEditsPPInversionPipelineOutput
[[autodoc]] pipelines.ledits_pp.pipeline_output.LEditsPPInversionPipelineOutput
- all

View File

@@ -57,7 +57,6 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Latent Consistency Models](latent_consistency_models) | text2image |
| [Latent Diffusion](latent_diffusion) | text2image, super-resolution |
| [LDM3D](stable_diffusion/ldm3d_diffusion) | text2image, text-to-3D, text-to-pano, upscaling |
| [LEDITS++](ledits_pp) | image editing |
| [MultiDiffusion](panorama) | text2image |
| [MusicLDM](musicldm) | text2audio |
| [Paint by Example](paint_by_example) | inpainting |
@@ -97,11 +96,6 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
- to
- components
[[autodoc]] pipelines.StableDiffusionMixin.enable_freeu
[[autodoc]] pipelines.StableDiffusionMixin.disable_freeu
## FlaxDiffusionPipeline
[[autodoc]] pipelines.pipeline_flax_utils.FlaxDiffusionPipeline

View File

@@ -30,6 +30,6 @@ Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers)
- all
- __call__
## SemanticStableDiffusionPipelineOutput
## StableDiffusionSafePipelineOutput
[[autodoc]] pipelines.semantic_stable_diffusion.pipeline_output.SemanticStableDiffusionPipelineOutput
- all

View File

@@ -1,229 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Stable Cascade
This model is built upon the [Würstchen](https://openreview.net/forum?id=gU58d5QeGv) architecture and its main
difference to other models like Stable Diffusion is that it is working at a much smaller latent space. Why is this
important? The smaller the latent space, the **faster** you can run inference and the **cheaper** the training becomes.
How small is the latent space? Stable Diffusion uses a compression factor of 8, resulting in a 1024x1024 image being
encoded to 128x128. Stable Cascade achieves a compression factor of 42, meaning that it is possible to encode a
1024x1024 image to 24x24, while maintaining crisp reconstructions. The text-conditional model is then trained in the
highly compressed latent space. Previous versions of this architecture, achieved a 16x cost reduction over Stable
Diffusion 1.5.
Therefore, this kind of model is well suited for usages where efficiency is important. Furthermore, all known extensions
like finetuning, LoRA, ControlNet, IP-Adapter, LCM etc. are possible with this method as well.
The original codebase can be found at [Stability-AI/StableCascade](https://github.com/Stability-AI/StableCascade).
## Model Overview
Stable Cascade consists of three models: Stage A, Stage B and Stage C, representing a cascade to generate images,
hence the name "Stable Cascade".
Stage A & B are used to compress images, similar to what the job of the VAE is in Stable Diffusion.
However, with this setup, a much higher compression of images can be achieved. While the Stable Diffusion models use a
spatial compression factor of 8, encoding an image with resolution of 1024 x 1024 to 128 x 128, Stable Cascade achieves
a compression factor of 42. This encodes a 1024 x 1024 image to 24 x 24, while being able to accurately decode the
image. This comes with the great benefit of cheaper training and inference. Furthermore, Stage C is responsible
for generating the small 24 x 24 latents given a text prompt.
The Stage C model operates on the small 24 x 24 latents and denoises the latents conditioned on text prompts. The model is also the largest component in the Cascade pipeline and is meant to be used with the `StableCascadePriorPipeline`
The Stage B and Stage A models are used with the `StableCascadeDecoderPipeline` and are responsible for generating the final image given the small 24 x 24 latents.
<Tip warning={true}>
There are some restrictions on data types that can be used with the Stable Cascade models. The official checkpoints for the `StableCascadePriorPipeline` do not support the `torch.float16` data type. Please use `torch.bfloat16` instead.
In order to use the `torch.bfloat16` data type with the `StableCascadeDecoderPipeline` you need to have PyTorch 2.2.0 or higher installed. This also means that using the `StableCascadeCombinedPipeline` with `torch.bfloat16` requires PyTorch 2.2.0 or higher, since it calls the `StableCascadeDecoderPipeline` internally.
If it is not possible to install PyTorch 2.2.0 or higher in your environment, the `StableCascadeDecoderPipeline` can be used on its own with the `torch.float16` data type. You can download the full precision or `bf16` variant weights for the pipeline and cast the weights to `torch.float16`.
</Tip>
## Usage example
```python
import torch
from diffusers import StableCascadeDecoderPipeline, StableCascadePriorPipeline
prompt = "an image of a shiba inu, donning a spacesuit and helmet"
negative_prompt = ""
prior = StableCascadePriorPipeline.from_pretrained("stabilityai/stable-cascade-prior", variant="bf16", torch_dtype=torch.bfloat16)
decoder = StableCascadeDecoderPipeline.from_pretrained("stabilityai/stable-cascade", variant="bf16", torch_dtype=torch.float16)
prior.enable_model_cpu_offload()
prior_output = prior(
prompt=prompt,
height=1024,
width=1024,
negative_prompt=negative_prompt,
guidance_scale=4.0,
num_images_per_prompt=1,
num_inference_steps=20
)
decoder.enable_model_cpu_offload()
decoder_output = decoder(
image_embeddings=prior_output.image_embeddings.to(torch.float16),
prompt=prompt,
negative_prompt=negative_prompt,
guidance_scale=0.0,
output_type="pil",
num_inference_steps=10
).images[0]
decoder_output.save("cascade.png")
```
## Using the Lite Versions of the Stage B and Stage C models
```python
import torch
from diffusers import (
StableCascadeDecoderPipeline,
StableCascadePriorPipeline,
StableCascadeUNet,
)
prompt = "an image of a shiba inu, donning a spacesuit and helmet"
negative_prompt = ""
prior_unet = StableCascadeUNet.from_pretrained("stabilityai/stable-cascade-prior", subfolder="prior_lite")
decoder_unet = StableCascadeUNet.from_pretrained("stabilityai/stable-cascade", subfolder="decoder_lite")
prior = StableCascadePriorPipeline.from_pretrained("stabilityai/stable-cascade-prior", prior=prior_unet)
decoder = StableCascadeDecoderPipeline.from_pretrained("stabilityai/stable-cascade", decoder=decoder_unet)
prior.enable_model_cpu_offload()
prior_output = prior(
prompt=prompt,
height=1024,
width=1024,
negative_prompt=negative_prompt,
guidance_scale=4.0,
num_images_per_prompt=1,
num_inference_steps=20
)
decoder.enable_model_cpu_offload()
decoder_output = decoder(
image_embeddings=prior_output.image_embeddings,
prompt=prompt,
negative_prompt=negative_prompt,
guidance_scale=0.0,
output_type="pil",
num_inference_steps=10
).images[0]
decoder_output.save("cascade.png")
```
## Loading original checkpoints with `from_single_file`
Loading the original format checkpoints is supported via `from_single_file` method in the StableCascadeUNet.
```python
import torch
from diffusers import (
StableCascadeDecoderPipeline,
StableCascadePriorPipeline,
StableCascadeUNet,
)
prompt = "an image of a shiba inu, donning a spacesuit and helmet"
negative_prompt = ""
prior_unet = StableCascadeUNet.from_single_file(
"https://huggingface.co/stabilityai/stable-cascade/resolve/main/stage_c_bf16.safetensors",
torch_dtype=torch.bfloat16
)
decoder_unet = StableCascadeUNet.from_single_file(
"https://huggingface.co/stabilityai/stable-cascade/blob/main/stage_b_bf16.safetensors",
torch_dtype=torch.bfloat16
)
prior = StableCascadePriorPipeline.from_pretrained("stabilityai/stable-cascade-prior", prior=prior_unet, torch_dtype=torch.bfloat16)
decoder = StableCascadeDecoderPipeline.from_pretrained("stabilityai/stable-cascade", decoder=decoder_unet, torch_dtype=torch.bfloat16)
prior.enable_model_cpu_offload()
prior_output = prior(
prompt=prompt,
height=1024,
width=1024,
negative_prompt=negative_prompt,
guidance_scale=4.0,
num_images_per_prompt=1,
num_inference_steps=20
)
decoder.enable_model_cpu_offload()
decoder_output = decoder(
image_embeddings=prior_output.image_embeddings,
prompt=prompt,
negative_prompt=negative_prompt,
guidance_scale=0.0,
output_type="pil",
num_inference_steps=10
).images[0]
decoder_output.save("cascade-single-file.png")
```
## Uses
### Direct Use
The model is intended for research purposes for now. Possible research areas and tasks include
- Research on generative models.
- Safe deployment of models which have the potential to generate harmful content.
- Probing and understanding the limitations and biases of generative models.
- Generation of artworks and use in design and other artistic processes.
- Applications in educational or creative tools.
Excluded uses are described below.
### Out-of-Scope Use
The model was not trained to be factual or true representations of people or events,
and therefore using the model to generate such content is out-of-scope for the abilities of this model.
The model should not be used in any way that violates Stability AI's [Acceptable Use Policy](https://stability.ai/use-policy).
## Limitations and Bias
### Limitations
- Faces and people in general may not be generated properly.
- The autoencoding part of the model is lossy.
## StableCascadeCombinedPipeline
[[autodoc]] StableCascadeCombinedPipeline
- all
- __call__
## StableCascadePriorPipeline
[[autodoc]] StableCascadePriorPipeline
- all
- __call__
## StableCascadePriorPipelineOutput
[[autodoc]] pipelines.stable_cascade.pipeline_stable_cascade_prior.StableCascadePriorPipelineOutput
## StableCascadeDecoderPipeline
[[autodoc]] StableCascadeDecoderPipeline
- all
- __call__

View File

@@ -10,7 +10,9 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# T2I-Adapter
# Text-to-Image Generation with Adapter Conditioning
## Overview
[T2I-Adapter: Learning Adapters to Dig out More Controllable Ability for Text-to-Image Diffusion Models](https://arxiv.org/abs/2302.08453) by Chong Mou, Xintao Wang, Liangbin Xie, Jian Zhang, Zhongang Qi, Ying Shan, Xiaohu Qie.
@@ -22,26 +24,236 @@ The abstract of the paper is the following:
This model was contributed by the community contributor [HimariO](https://github.com/HimariO) ❤️ .
## StableDiffusionAdapterPipeline
## Available Pipelines:
| Pipeline | Tasks | Demo
|---|---|:---:|
| [StableDiffusionAdapterPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_adapter.py) | *Text-to-Image Generation with T2I-Adapter Conditioning* | -
| [StableDiffusionXLAdapterPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_xl_adapter.py) | *Text-to-Image Generation with T2I-Adapter Conditioning on StableDiffusion-XL* | -
## Usage example with the base model of StableDiffusion-1.4/1.5
In the following we give a simple example of how to use a *T2I-Adapter* checkpoint with Diffusers for inference based on StableDiffusion-1.4/1.5.
All adapters use the same pipeline.
1. Images are first converted into the appropriate *control image* format.
2. The *control image* and *prompt* are passed to the [`StableDiffusionAdapterPipeline`].
Let's have a look at a simple example using the [Color Adapter](https://huggingface.co/TencentARC/t2iadapter_color_sd14v1).
```python
from diffusers.utils import load_image, make_image_grid
image = load_image("https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_ref.png")
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_ref.png)
Then we can create our color palette by simply resizing it to 8 by 8 pixels and then scaling it back to original size.
```python
from PIL import Image
color_palette = image.resize((8, 8))
color_palette = color_palette.resize((512, 512), resample=Image.Resampling.NEAREST)
```
Let's take a look at the processed image.
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_palette.png)
Next, create the adapter pipeline
```py
import torch
from diffusers import StableDiffusionAdapterPipeline, T2IAdapter
adapter = T2IAdapter.from_pretrained("TencentARC/t2iadapter_color_sd14v1", torch_dtype=torch.float16)
pipe = StableDiffusionAdapterPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
adapter=adapter,
torch_dtype=torch.float16,
)
pipe.to("cuda")
```
Finally, pass the prompt and control image to the pipeline
```py
# fix the random seed, so you will get the same result as the example
generator = torch.Generator("cuda").manual_seed(7)
out_image = pipe(
"At night, glowing cubes in front of the beach",
image=color_palette,
generator=generator,
).images[0]
make_image_grid([image, color_palette, out_image], rows=1, cols=3)
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_output.png)
## Usage example with the base model of StableDiffusion-XL
In the following we give a simple example of how to use a *T2I-Adapter* checkpoint with Diffusers for inference based on StableDiffusion-XL.
All adapters use the same pipeline.
1. Images are first downloaded into the appropriate *control image* format.
2. The *control image* and *prompt* are passed to the [`StableDiffusionXLAdapterPipeline`].
Let's have a look at a simple example using the [Sketch Adapter](https://huggingface.co/Adapter/t2iadapter/tree/main/sketch_sdxl_1.0).
```python
from diffusers.utils import load_image, make_image_grid
sketch_image = load_image("https://huggingface.co/Adapter/t2iadapter/resolve/main/sketch.png").convert("L")
```
![img](https://huggingface.co/Adapter/t2iadapter/resolve/main/sketch.png)
Then, create the adapter pipeline
```py
import torch
from diffusers import (
T2IAdapter,
StableDiffusionXLAdapterPipeline,
DDPMScheduler
)
model_id = "stabilityai/stable-diffusion-xl-base-1.0"
adapter = T2IAdapter.from_pretrained("Adapter/t2iadapter", subfolder="sketch_sdxl_1.0", torch_dtype=torch.float16, adapter_type="full_adapter_xl")
scheduler = DDPMScheduler.from_pretrained(model_id, subfolder="scheduler")
pipe = StableDiffusionXLAdapterPipeline.from_pretrained(
model_id, adapter=adapter, safety_checker=None, torch_dtype=torch.float16, variant="fp16", scheduler=scheduler
)
pipe.to("cuda")
```
Finally, pass the prompt and control image to the pipeline
```py
# fix the random seed, so you will get the same result as the example
generator = torch.Generator().manual_seed(42)
sketch_image_out = pipe(
prompt="a photo of a dog in real world, high quality",
negative_prompt="extra digit, fewer digits, cropped, worst quality, low quality",
image=sketch_image,
generator=generator,
guidance_scale=7.5
).images[0]
make_image_grid([sketch_image, sketch_image_out], rows=1, cols=2)
```
![img](https://huggingface.co/Adapter/t2iadapter/resolve/main/sketch_output.png)
## Available checkpoints
Non-diffusers checkpoints can be found under [TencentARC/T2I-Adapter](https://huggingface.co/TencentARC/T2I-Adapter/tree/main/models).
### T2I-Adapter with Stable Diffusion 1.4
| Model Name | Control Image Overview| Control Image Example | Generated Image Example |
|---|---|---|---|
|[TencentARC/t2iadapter_color_sd14v1](https://huggingface.co/TencentARC/t2iadapter_color_sd14v1)<br/> *Trained with spatial color palette* | An image with 8x8 color palette.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_sample_input.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/color_sample_output.png"/></a>|
|[TencentARC/t2iadapter_canny_sd14v1](https://huggingface.co/TencentARC/t2iadapter_canny_sd14v1)<br/> *Trained with canny edge detection* | A monochrome image with white edges on a black background.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/canny_sample_input.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/canny_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/canny_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/canny_sample_output.png"/></a>|
|[TencentARC/t2iadapter_sketch_sd14v1](https://huggingface.co/TencentARC/t2iadapter_sketch_sd14v1)<br/> *Trained with [PidiNet](https://github.com/zhuoinoulu/pidinet) edge detection* | A hand-drawn monochrome image with white outlines on a black background.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/sketch_sample_input.png"><img width="64" style="margin:0;padding:0;" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/sketch_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/sketch_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/sketch_sample_output.png"/></a>|
|[TencentARC/t2iadapter_depth_sd14v1](https://huggingface.co/TencentARC/t2iadapter_depth_sd14v1)<br/> *Trained with Midas depth estimation* | A grayscale image with black representing deep areas and white representing shallow areas.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_input.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_output.png"/></a>|
|[TencentARC/t2iadapter_openpose_sd14v1](https://huggingface.co/TencentARC/t2iadapter_openpose_sd14v1)<br/> *Trained with OpenPose bone image* | A [OpenPose bone](https://github.com/CMU-Perceptual-Computing-Lab/openpose) image.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/openpose_sample_input.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/openpose_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/openpose_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/openpose_sample_output.png"/></a>|
|[TencentARC/t2iadapter_keypose_sd14v1](https://huggingface.co/TencentARC/t2iadapter_keypose_sd14v1)<br/> *Trained with mmpose skeleton image* | A [mmpose skeleton](https://github.com/open-mmlab/mmpose) image.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_input.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_output.png"/></a>|
|[TencentARC/t2iadapter_seg_sd14v1](https://huggingface.co/TencentARC/t2iadapter_seg_sd14v1)<br/>*Trained with semantic segmentation* | An [custom](https://github.com/TencentARC/T2I-Adapter/discussions/25) segmentation protocol image.|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/seg_sample_input.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/seg_sample_input.png"/></a>|<a href="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/seg_sample_output.png"><img width="64" src="https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/seg_sample_output.png"/></a> |
|[TencentARC/t2iadapter_canny_sd15v2](https://huggingface.co/TencentARC/t2iadapter_canny_sd15v2)||
|[TencentARC/t2iadapter_depth_sd15v2](https://huggingface.co/TencentARC/t2iadapter_depth_sd15v2)||
|[TencentARC/t2iadapter_sketch_sd15v2](https://huggingface.co/TencentARC/t2iadapter_sketch_sd15v2)||
|[TencentARC/t2iadapter_zoedepth_sd15v1](https://huggingface.co/TencentARC/t2iadapter_zoedepth_sd15v1)||
|[Adapter/t2iadapter, subfolder='sketch_sdxl_1.0'](https://huggingface.co/Adapter/t2iadapter/tree/main/sketch_sdxl_1.0)||
|[Adapter/t2iadapter, subfolder='canny_sdxl_1.0'](https://huggingface.co/Adapter/t2iadapter/tree/main/canny_sdxl_1.0)||
|[Adapter/t2iadapter, subfolder='openpose_sdxl_1.0'](https://huggingface.co/Adapter/t2iadapter/tree/main/openpose_sdxl_1.0)||
## Combining multiple adapters
[`MultiAdapter`] can be used for applying multiple conditionings at once.
Here we use the keypose adapter for the character posture and the depth adapter for creating the scene.
```py
from diffusers.utils import load_image, make_image_grid
cond_keypose = load_image(
"https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_input.png"
)
cond_depth = load_image(
"https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_input.png"
)
cond = [cond_keypose, cond_depth]
prompt = ["A man walking in an office room with a nice view"]
```
The two control images look as such:
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_sample_input.png)
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/depth_sample_input.png)
`MultiAdapter` combines keypose and depth adapters.
`adapter_conditioning_scale` balances the relative influence of the different adapters.
```py
import torch
from diffusers import StableDiffusionAdapterPipeline, MultiAdapter, T2IAdapter
adapters = MultiAdapter(
[
T2IAdapter.from_pretrained("TencentARC/t2iadapter_keypose_sd14v1"),
T2IAdapter.from_pretrained("TencentARC/t2iadapter_depth_sd14v1"),
]
)
adapters = adapters.to(torch.float16)
pipe = StableDiffusionAdapterPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
torch_dtype=torch.float16,
adapter=adapters,
).to("cuda")
image = pipe(prompt, cond, adapter_conditioning_scale=[0.8, 0.8]).images[0]
make_image_grid([cond_keypose, cond_depth, image], rows=1, cols=3)
```
![img](https://huggingface.co/datasets/diffusers/docs-images/resolve/main/t2i-adapter/keypose_depth_sample_output.png)
## T2I-Adapter vs ControlNet
T2I-Adapter is similar to [ControlNet](https://huggingface.co/docs/diffusers/main/en/api/pipelines/controlnet).
T2I-Adapter uses a smaller auxiliary network which is only run once for the entire diffusion process.
However, T2I-Adapter performs slightly worse than ControlNet.
## StableDiffusionAdapterPipeline
[[autodoc]] StableDiffusionAdapterPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_vae_slicing
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_vae_slicing
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
## StableDiffusionXLAdapterPipeline
[[autodoc]] StableDiffusionXLAdapterPipeline
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_vae_slicing
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention
- all
- __call__
- enable_attention_slicing
- disable_attention_slicing
- enable_vae_slicing
- disable_vae_slicing
- enable_xformers_memory_efficient_attention
- disable_xformers_memory_efficient_attention

View File

@@ -172,41 +172,3 @@ inpaint = StableDiffusionInpaintPipeline(**text2img.components)
# now you can use text2img(...), img2img(...), inpaint(...) just like the call methods of each respective pipeline
```
### Create web demos using `gradio`
The Stable Diffusion pipelines are automatically supported in [Gradio](https://github.com/gradio-app/gradio/), a library that makes creating beautiful and user-friendly machine learning apps on the web a breeze. First, make sure you have Gradio installed:
```
pip install -U gradio
```
Then, create a web demo around any Stable Diffusion-based pipeline. For example, you can create an image generation pipeline in a single line of code with Gradio's [`Interface.from_pipeline`](https://www.gradio.app/docs/interface#interface-from-pipeline) function:
```py
from diffusers import StableDiffusionPipeline
import gradio as gr
pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4")
gr.Interface.from_pipeline(pipe).launch()
```
which opens an intuitive drag-and-drop interface in your browser:
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/gradio-panda.png)
Similarly, you could create a demo for an image-to-image pipeline with:
```py
from diffusers import StableDiffusionImg2ImgPipeline
import gradio as gr
pipe = StableDiffusionImg2ImgPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
gr.Interface.from_pipeline(pipe).launch()
```
By default, the web demo runs on a local server. If you'd like to share it with others, you can generate a temporary public
link by setting `share=True` in `launch()`. Or, you can host your demo on [Hugging Face Spaces](https://huggingface.co/spaces)https://huggingface.co/spaces for a permanent link.

View File

@@ -21,7 +21,7 @@ The abstract from the paper is:
## Tips
- SDXL Turbo uses the exact same architecture as [SDXL](./stable_diffusion_xl), which means it also has the same API. Please refer to the [SDXL](./stable_diffusion_xl) API reference for more details.
- SDXL Turbo should disable guidance scale by setting `guidance_scale=0.0`.
- SDXL Turbo should disable guidance scale by setting `guidance_scale=0.0`
- SDXL Turbo should use `timestep_spacing='trailing'` for the scheduler and use between 1 and 4 steps.
- SDXL Turbo has been trained to generate images of size 512x512.
- SDXL Turbo is open-access, but not open-source meaning that one might have to buy a model license in order to use it for commercial applications. Make sure to read the [official model card](https://huggingface.co/stabilityai/sdxl-turbo) to learn more.

View File

@@ -1,21 +1,9 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ConsistencyDecoderScheduler
This scheduler is a part of the [`ConsistencyDecoderPipeline`] and was introduced in [DALL-E 3](https://openai.com/dall-e-3).
This scheduler is a part of the [`ConsistencyDecoderPipeline`] and was introduced in [DALL-E 3](https://openai.com/dall-e-3).
The original codebase can be found at [openai/consistency_models](https://github.com/openai/consistency_models).
## ConsistencyDecoderScheduler
[[autodoc]] schedulers.scheduling_consistency_decoder.ConsistencyDecoderScheduler
[[autodoc]] schedulers.scheduling_consistency_decoder.ConsistencyDecoderScheduler

View File

@@ -1,22 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# EDMEulerScheduler
The Karras formulation of the Euler scheduler (Algorithm 2) from the [Elucidating the Design Space of Diffusion-Based Generative Models](https://huggingface.co/papers/2206.00364) paper by Karras et al. This is a fast scheduler which can often generate good outputs in 20-30 steps. The scheduler is based on the original [k-diffusion](https://github.com/crowsonkb/k-diffusion/blob/481677d114f6ea445aa009cf5bd7a9cdee909e47/k_diffusion/sampling.py#L51) implementation by [Katherine Crowson](https://github.com/crowsonkb/).
## EDMEulerScheduler
[[autodoc]] EDMEulerScheduler
## EDMEulerSchedulerOutput
[[autodoc]] schedulers.scheduling_edm_euler.EDMEulerSchedulerOutput

View File

@@ -1,24 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# EDMDPMSolverMultistepScheduler
`EDMDPMSolverMultistepScheduler` is a [Karras formulation](https://huggingface.co/papers/2206.00364) of `DPMSolverMultistep`, a multistep scheduler from [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095) by Cheng Lu, Yuhao Zhou, Fan Bao, Jianfei Chen, Chongxuan Li, and Jun Zhu.
DPMSolver (and the improved version DPMSolver++) is a fast dedicated high-order solver for diffusion ODEs with convergence order guarantee. Empirically, DPMSolver sampling with only 20 steps can generate high-quality
samples, and it can generate quite good samples even in 10 steps.
## EDMDPMSolverMultistepScheduler
[[autodoc]] EDMDPMSolverMultistepScheduler
## SchedulerOutput
[[autodoc]] schedulers.scheduling_utils.SchedulerOutput

View File

@@ -1,29 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# TCDScheduler
[Trajectory Consistency Distillation](https://huggingface.co/papers/2402.19159) by Jianbin Zheng, Minghui Hu, Zhongyi Fan, Chaoyue Wang, Changxing Ding, Dacheng Tao and Tat-Jen Cham introduced a Strategic Stochastic Sampling (Algorithm 4) that is capable of generating good samples in a small number of steps. Distinguishing it as an advanced iteration of the multistep scheduler (Algorithm 1) in the [Consistency Models](https://huggingface.co/papers/2303.01469), Strategic Stochastic Sampling specifically tailored for the trajectory consistency function.
The abstract from the paper is:
*Latent Consistency Model (LCM) extends the Consistency Model to the latent space and leverages the guided consistency distillation technique to achieve impressive performance in accelerating text-to-image synthesis. However, we observed that LCM struggles to generate images with both clarity and detailed intricacy. To address this limitation, we initially delve into and elucidate the underlying causes. Our investigation identifies that the primary issue stems from errors in three distinct areas. Consequently, we introduce Trajectory Consistency Distillation (TCD), which encompasses trajectory consistency function and strategic stochastic sampling. The trajectory consistency function diminishes the distillation errors by broadening the scope of the self-consistency boundary condition and endowing the TCD with the ability to accurately trace the entire trajectory of the Probability Flow ODE. Additionally, strategic stochastic sampling is specifically designed to circumvent the accumulated errors inherent in multi-step consistency sampling, which is meticulously tailored to complement the TCD model. Experiments demonstrate that TCD not only significantly enhances image quality at low NFEs but also yields more detailed results compared to the teacher model at high NFEs.*
The original codebase can be found at [jabir-zheng/TCD](https://github.com/jabir-zheng/TCD).
## TCDScheduler
[[autodoc]] TCDScheduler
## TCDSchedulerOutput
[[autodoc]] schedulers.scheduling_tcd.TCDSchedulerOutput

View File

@@ -37,7 +37,3 @@ Utility and helper functions for working with 🤗 Diffusers.
## make_image_grid
[[autodoc]] utils.make_image_grid
## randn_tensor
[[autodoc]] utils.torch_utils.randn_tensor

View File

@@ -198,81 +198,38 @@ Anything displayed on [the official Diffusers doc page](https://huggingface.co/d
Please have a look at [this page](https://github.com/huggingface/diffusers/tree/main/docs) on how to verify changes made to the documentation locally.
### 6. Contribute a community pipeline
> [!TIP]
> Read the [Community pipelines](../using-diffusers/custom_pipeline_overview#community-pipelines) guide to learn more about the difference between a GitHub and Hugging Face Hub community pipeline. If you're interested in why we have community pipelines, take a look at GitHub Issue [#841](https://github.com/huggingface/diffusers/issues/841) (basically, we can't maintain all the possible ways diffusion models can be used for inference but we also don't want to prevent the community from building them).
[Pipelines](https://huggingface.co/docs/diffusers/api/pipelines/overview) are usually the first point of contact between the Diffusers library and the user.
Pipelines are examples of how to use Diffusers [models](https://huggingface.co/docs/diffusers/api/models/overview) and [schedulers](https://huggingface.co/docs/diffusers/api/schedulers/overview).
We support two types of pipelines:
Contributing a community pipeline is a great way to share your creativity and work with the community. It lets you build on top of the [`DiffusionPipeline`] so that anyone can load and use it by setting the `custom_pipeline` parameter. This section will walk you through how to create a simple pipeline where the UNet only does a single forward pass and calls the scheduler once (a "one-step" pipeline).
- Official Pipelines
- Community Pipelines
1. Create a one_step_unet.py file for your community pipeline. This file can contain whatever package you want to use as long as it's installed by the user. Make sure you only have one pipeline class that inherits from [`DiffusionPipeline`] to load model weights and the scheduler configuration from the Hub. Add a UNet and scheduler to the `__init__` function.
Both official and community pipelines follow the same design and consist of the same type of components.
You should also add the `register_modules` function to ensure your pipeline and its components can be saved with [`~DiffusionPipeline.save_pretrained`].
Official pipelines are tested and maintained by the core maintainers of Diffusers. Their code
resides in [src/diffusers/pipelines](https://github.com/huggingface/diffusers/tree/main/src/diffusers/pipelines).
In contrast, community pipelines are contributed and maintained purely by the **community** and are **not** tested.
They reside in [examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community) and while they can be accessed via the [PyPI diffusers package](https://pypi.org/project/diffusers/), their code is not part of the PyPI distribution.
```py
from diffusers import DiffusionPipeline
import torch
The reason for the distinction is that the core maintainers of the Diffusers library cannot maintain and test all
possible ways diffusion models can be used for inference, but some of them may be of interest to the community.
Officially released diffusion pipelines,
such as Stable Diffusion are added to the core src/diffusers/pipelines package which ensures
high quality of maintenance, no backward-breaking code changes, and testing.
More bleeding edge pipelines should be added as community pipelines. If usage for a community pipeline is high, the pipeline can be moved to the official pipelines upon request from the community. This is one of the ways we strive to be a community-driven library.
class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
def __init__(self, unet, scheduler):
super().__init__()
To add a community pipeline, one should add a <name-of-the-community>.py file to [examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community) and adapt the [examples/community/README.md](https://github.com/huggingface/diffusers/tree/main/examples/community/README.md) to include an example of the new pipeline.
self.register_modules(unet=unet, scheduler=scheduler)
```
An example can be seen [here](https://github.com/huggingface/diffusers/pull/2400).
1. In the forward pass (which we recommend defining as `__call__`), you can add any feature you'd like. For the "one-step" pipeline, create a random image and call the UNet and scheduler once by setting `timestep=1`.
Community pipeline PRs are only checked at a superficial level and ideally they should be maintained by their original authors.
```py
from diffusers import DiffusionPipeline
import torch
class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
def __init__(self, unet, scheduler):
super().__init__()
self.register_modules(unet=unet, scheduler=scheduler)
def __call__(self):
image = torch.randn(
(1, self.unet.config.in_channels, self.unet.config.sample_size, self.unet.config.sample_size),
)
timestep = 1
model_output = self.unet(image, timestep).sample
scheduler_output = self.scheduler.step(model_output, timestep, image).prev_sample
return scheduler_output
```
Now you can run the pipeline by passing a UNet and scheduler to it or load pretrained weights if the pipeline structure is identical.
```py
from diffusers import DDPMScheduler, UNet2DModel
scheduler = DDPMScheduler()
unet = UNet2DModel()
pipeline = UnetSchedulerOneForwardPipeline(unet=unet, scheduler=scheduler)
output = pipeline()
# load pretrained weights
pipeline = UnetSchedulerOneForwardPipeline.from_pretrained("google/ddpm-cifar10-32", use_safetensors=True)
output = pipeline()
```
You can either share your pipeline as a GitHub community pipeline or Hub community pipeline.
<hfoptions id="pipeline type">
<hfoption id="GitHub pipeline">
Share your GitHub pipeline by opening a pull request on the Diffusers [repository](https://github.com/huggingface/diffusers) and add the one_step_unet.py file to the [examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community) subfolder.
</hfoption>
<hfoption id="Hub pipeline">
Share your Hub pipeline by creating a model repository on the Hub and uploading the one_step_unet.py file to it.
</hfoption>
</hfoptions>
Contributing a community pipeline is a great way to understand how Diffusers models and schedulers work. Having contributed a community pipeline is usually the first stepping stone to contributing an official pipeline to the
core package.
### 7. Contribute to training examples

View File

@@ -66,9 +66,3 @@ image = pipe(prompt).images[0]
Don't use [`torch.autocast`](https://pytorch.org/docs/stable/amp.html#torch.autocast) in any of the pipelines as it can lead to black images and is always slower than pure float16 precision.
</Tip>
## Distilled model
You could also use a distilled Stable Diffusion model and autoencoder to speed up inference. During distillation, many of the UNet's residual and attention blocks are shed to reduce the model size. The distilled model is faster and uses less memory while generating images of comparable quality to the full Stable Diffusion model.
Learn more about in the [Distilled Stable Diffusion inference](../using-diffusers/distilled_sd) guide!

View File

@@ -0,0 +1,17 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Overview
Generating high-quality outputs is computationally intensive, especially during each iterative step where you go from a noisy output to a less noisy output. One of 🤗 Diffuser's goals is to make this technology widely accessible to everyone, which includes enabling fast inference on consumer and specialized hardware.
This section will cover tips and tricks - like half-precision weights and sliced attention - for optimizing inference speed and reducing memory-consumption. You'll also learn how to speed up your PyTorch code with [`torch.compile`](https://pytorch.org/tutorials/intermediate/torch_compile_tutorial.html) or [ONNX Runtime](https://onnxruntime.ai/docs/), and enable memory-efficient attention with [xFormers](https://facebookresearch.github.io/xformers/). There are also guides for running inference on specific hardware like Apple Silicon, and Intel or Habana processors.

View File

@@ -1,182 +0,0 @@
# T-GATE
[T-GATE](https://github.com/HaozheLiu-ST/T-GATE/tree/main) accelerates inference for [Stable Diffusion](../api/pipelines/stable_diffusion/overview), [PixArt](../api/pipelines/pixart), and [Latency Consistency Model](../api/pipelines/latent_consistency_models.md) pipelines by skipping the cross-attention calculation once it converges. This method doesn't require any additional training and it can speed up inference from 10-50%. T-GATE is also compatible with other optimization methods like [DeepCache](./deepcache).
Before you begin, make sure you install T-GATE.
```bash
pip install tgate
pip install -U pytorch diffusers transformers accelerate DeepCache
```
To use T-GATE with a pipeline, you need to use its corresponding loader.
| Pipeline | T-GATE Loader |
|---|---|
| PixArt | TgatePixArtLoader |
| Stable Diffusion XL | TgateSDXLLoader |
| Stable Diffusion XL + DeepCache | TgateSDXLDeepCacheLoader |
| Stable Diffusion | TgateSDLoader |
| Stable Diffusion + DeepCache | TgateSDDeepCacheLoader |
Next, create a `TgateLoader` with a pipeline, the gate step (the time step to stop calculating the cross attention), and the number of inference steps. Then call the `tgate` method on the pipeline with a prompt, gate step, and the number of inference steps.
Let's see how to enable this for several different pipelines.
<hfoptions id="pipelines">
<hfoption id="PixArt">
Accelerate `PixArtAlphaPipeline` with T-GATE:
```py
import torch
from diffusers import PixArtAlphaPipeline
from tgate import TgatePixArtLoader
pipe = PixArtAlphaPipeline.from_pretrained("PixArt-alpha/PixArt-XL-2-1024-MS", torch_dtype=torch.float16)
gate_step = 8
inference_step = 25
pipe = TgatePixArtLoader(
pipe,
gate_step=gate_step,
num_inference_steps=inference_step,
).to("cuda")
image = pipe.tgate(
"An alpaca made of colorful building blocks, cyberpunk.",
gate_step=gate_step,
num_inference_steps=inference_step,
).images[0]
```
</hfoption>
<hfoption id="Stable Diffusion XL">
Accelerate `StableDiffusionXLPipeline` with T-GATE:
```py
import torch
from diffusers import StableDiffusionXLPipeline
from diffusers import DPMSolverMultistepScheduler
from tgate import TgateSDXLLoader
pipe = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
variant="fp16",
use_safetensors=True,
)
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
gate_step = 10
inference_step = 25
pipe = TgateSDXLLoader(
pipe,
gate_step=gate_step,
num_inference_steps=inference_step,
).to("cuda")
image = pipe.tgate(
"Astronaut in a jungle, cold color palette, muted colors, detailed, 8k.",
gate_step=gate_step,
num_inference_steps=inference_step
).images[0]
```
</hfoption>
<hfoption id="StableDiffusionXL with DeepCache">
Accelerate `StableDiffusionXLPipeline` with [DeepCache](https://github.com/horseee/DeepCache) and T-GATE:
```py
import torch
from diffusers import StableDiffusionXLPipeline
from diffusers import DPMSolverMultistepScheduler
from tgate import TgateSDXLDeepCacheLoader
pipe = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
variant="fp16",
use_safetensors=True,
)
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
gate_step = 10
inference_step = 25
pipe = TgateSDXLDeepCacheLoader(
pipe,
cache_interval=3,
cache_branch_id=0,
).to("cuda")
image = pipe.tgate(
"Astronaut in a jungle, cold color palette, muted colors, detailed, 8k.",
gate_step=gate_step,
num_inference_steps=inference_step
).images[0]
```
</hfoption>
<hfoption id="Latent Consistency Model">
Accelerate `latent-consistency/lcm-sdxl` with T-GATE:
```py
import torch
from diffusers import StableDiffusionXLPipeline
from diffusers import UNet2DConditionModel, LCMScheduler
from diffusers import DPMSolverMultistepScheduler
from tgate import TgateSDXLLoader
unet = UNet2DConditionModel.from_pretrained(
"latent-consistency/lcm-sdxl",
torch_dtype=torch.float16,
variant="fp16",
)
pipe = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
unet=unet,
torch_dtype=torch.float16,
variant="fp16",
)
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
gate_step = 1
inference_step = 4
pipe = TgateSDXLLoader(
pipe,
gate_step=gate_step,
num_inference_steps=inference_step,
lcm=True
).to("cuda")
image = pipe.tgate(
"Astronaut in a jungle, cold color palette, muted colors, detailed, 8k.",
gate_step=gate_step,
num_inference_steps=inference_step
).images[0]
```
</hfoption>
</hfoptions>
T-GATE also supports [`StableDiffusionPipeline`] and [PixArt-alpha/PixArt-LCM-XL-2-1024-MS](https://hf.co/PixArt-alpha/PixArt-LCM-XL-2-1024-MS).
## Benchmarks
| Model | MACs | Param | Latency | Zero-shot 10K-FID on MS-COCO |
|-----------------------|----------|-----------|---------|---------------------------|
| SD-1.5 | 16.938T | 859.520M | 7.032s | 23.927 |
| SD-1.5 w/ T-GATE | 9.875T | 815.557M | 4.313s | 20.789 |
| SD-2.1 | 38.041T | 865.785M | 16.121s | 22.609 |
| SD-2.1 w/ T-GATE | 22.208T | 815.433 M | 9.878s | 19.940 |
| SD-XL | 149.438T | 2.570B | 53.187s | 24.628 |
| SD-XL w/ T-GATE | 84.438T | 2.024B | 27.932s | 22.738 |
| Pixart-Alpha | 107.031T | 611.350M | 61.502s | 38.669 |
| Pixart-Alpha w/ T-GATE | 65.318T | 462.585M | 37.867s | 35.825 |
| DeepCache (SD-XL) | 57.888T | - | 19.931s | 23.755 |
| DeepCache w/ T-GATE | 43.868T | - | 14.666s | 23.999 |
| LCM (SD-XL) | 11.955T | 2.570B | 3.805s | 25.044 |
| LCM w/ T-GATE | 11.171T | 2.024B | 3.533s | 25.028 |
| LCM (Pixart-Alpha) | 8.563T | 611.350M | 4.733s | 36.086 |
| LCM w/ T-GATE | 7.623T | 462.585M | 4.543s | 37.048 |
The latency is tested on an NVIDIA 1080TI, MACs and Params are calculated with [calflops](https://github.com/MrYxJ/calculate-flops.pytorch), and the FID is calculated with [PytorchFID](https://github.com/mseitzer/pytorch-fid).

View File

@@ -75,9 +75,6 @@ Compilation requires some time to complete, so it is best suited for situations
For more information and different options about `torch.compile`, refer to the [`torch_compile`](https://pytorch.org/tutorials/intermediate/torch_compile_tutorial.html) tutorial.
> [!TIP]
> Learn more about other ways PyTorch 2.0 can help optimize your model in the [Accelerate inference of text-to-image diffusion models](../tutorials/fast_diffusion) tutorial.
## Benchmark
We conducted a comprehensive benchmark with PyTorch 2.0's efficient attention implementation and `torch.compile` across different GPUs and batch sizes for five of our most used pipelines. The code is benchmarked on 🤗 Diffusers v0.17.0.dev0 to optimize `torch.compile` usage (see [here](https://github.com/huggingface/diffusers/pull/3313) for more details).

View File

@@ -49,7 +49,7 @@ One of the simplest ways to speed up inference is to place the pipeline on a GPU
pipeline = pipeline.to("cuda")
```
To make sure you can use the same image and improve on it, use a [`Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) and set a seed for [reproducibility](./using-diffusers/reusing_seeds):
To make sure you can use the same image and improve on it, use a [`Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) and set a seed for [reproducibility](./using-diffusers/reproducibility):
```python
import torch

View File

@@ -88,7 +88,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()

View File

@@ -54,7 +54,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -84,7 +84,7 @@ Many of the basic parameters are described in the [DreamBooth](dreambooth#script
- `--freeze_model`: freezes the key and value parameters in the cross-attention layer; the default is `crossattn_kv`, but you can set it to `crossattn` to train all the parameters in the cross-attention layer
- `--concepts_list`: to learn multiple concepts, provide a path to a JSON file containing the concepts
- `--modifier_token`: a special word used to represent the learned concept
- `--initializer_token`: a special word used to initialize the embeddings of the `modifier_token`
- `--initializer_token`:
### Prior preservation loss

View File

@@ -52,76 +52,6 @@ To learn more, take a look at the [Distributed Inference with 🤗 Accelerate](h
</Tip>
### Device placement
> [!WARNING]
> This feature is experimental and its APIs might change in the future.
With Accelerate, you can use the `device_map` to determine how to distribute the models of a pipeline across multiple devices. This is useful in situations where you have more than one GPU.
For example, if you have two 8GB GPUs, then using [`~DiffusionPipeline.enable_model_cpu_offload`] may not work so well because:
* it only works on a single GPU
* a single model might not fit on a single GPU ([`~DiffusionPipeline.enable_sequential_cpu_offload`] might work but it will be extremely slow and it is also limited to a single GPU)
To make use of both GPUs, you can use the "balanced" device placement strategy which splits the models across all available GPUs.
> [!WARNING]
> Only the "balanced" strategy is supported at the moment, and we plan to support additional mapping strategies in the future.
```diff
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained(
- "runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True,
+ "runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True, device_map="balanced"
)
image = pipeline("a dog").images[0]
image
```
You can also pass a dictionary to enforce the maximum GPU memory that can be used on each device:
```diff
from diffusers import DiffusionPipeline
import torch
max_memory = {0:"1GB", 1:"1GB"}
pipeline = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
device_map="balanced",
+ max_memory=max_memory
)
image = pipeline("a dog").images[0]
image
```
If a device is not present in `max_memory`, then it will be completely ignored and will not participate in the device placement.
By default, Diffusers uses the maximum memory of all devices. If the models don't fit on the GPUs, they are offloaded to the CPU. If the CPU doesn't have enough memory, then you might see an error. In that case, you could defer to using [`~DiffusionPipeline.enable_sequential_cpu_offload`] and [`~DiffusionPipeline.enable_model_cpu_offload`].
Call [`~DiffusionPipeline.reset_device_map`] to reset the `device_map` of a pipeline. This is also necessary if you want to use methods like `to()`, [`~DiffusionPipeline.enable_sequential_cpu_offload`], and [`~DiffusionPipeline.enable_model_cpu_offload`] on a pipeline that was device-mapped.
```py
pipeline.reset_device_map()
```
Once a pipeline has been device-mapped, you can also access its device map via `hf_device_map`:
```py
print(pipeline.hf_device_map)
```
An example device map would look like so:
```bash
{'unet': 1, 'vae': 1, 'safety_checker': 0, 'text_encoder': 0}
```
## PyTorch Distributed
PyTorch supports [`DistributedDataParallel`](https://pytorch.org/docs/stable/generated/torch.nn.parallel.DistributedDataParallel.html) which enables data parallelism.

View File

@@ -67,7 +67,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -180,7 +180,7 @@ elif args.pretrained_model_name_or_path:
revision=args.revision,
use_fast=False,
)
# Load scheduler and models
noise_scheduler = DDPMScheduler.from_pretrained(args.pretrained_model_name_or_path, subfolder="scheduler")
text_encoder = text_encoder_cls.from_pretrained(

View File

@@ -51,7 +51,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -89,7 +89,7 @@ The dataset preprocessing code and training loop are found in the [`main()`](htt
As with the script parameters, a walkthrough of the training script is provided in the [Text-to-image](text2image#training-script) training guide. Instead, this guide takes a look at the InstructPix2Pix relevant parts of the script.
The script begins by modifying the [number of input channels](https://github.com/huggingface/diffusers/blob/64603389da01082055a901f2883c4810d1144edb/examples/instruct_pix2pix/train_instruct_pix2pix.py#L445) in the first convolutional layer of the UNet to account for InstructPix2Pix's additional conditioning image:
The script begins by modifing the [number of input channels](https://github.com/huggingface/diffusers/blob/64603389da01082055a901f2883c4810d1144edb/examples/instruct_pix2pix/train_instruct_pix2pix.py#L445) in the first convolutional layer of the UNet to account for InstructPix2Pix's additional conditioning image:
```py
in_channels = 8

View File

@@ -59,7 +59,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -235,7 +235,7 @@ accelerate launch --mixed_precision="fp16" train_text_to_image_prior.py \
--validation_prompts="A robot pokemon, 4k photo" \
--report_to="wandb" \
--push_to_hub \
--output_dir="kandi2-prior-pokemon-model"
--output_dir="kandi2-prior-pokemon-model"
```
</hfoption>
@@ -259,7 +259,7 @@ accelerate launch --mixed_precision="fp16" train_text_to_image_decoder.py \
--validation_prompts="A robot pokemon, 4k photo" \
--report_to="wandb" \
--push_to_hub \
--output_dir="kandi2-decoder-pokemon-model"
--output_dir="kandi2-decoder-pokemon-model"
```
</hfoption>

View File

@@ -53,7 +53,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -252,4 +252,4 @@ The SDXL training script is discussed in more detail in the [SDXL training](sdxl
Congratulations on distilling a LCM model! To learn more about LCM, the following may be helpful:
- Learn how to use [LCMs for inference](../using-diffusers/lcm) for text-to-image, image-to-image, and with LoRA checkpoints.
- Read the [SDXL in 4 steps with Latent Consistency LoRAs](https://huggingface.co/blog/lcm_lora) blog post to learn more about SDXL LCM-LoRA's for super fast inference, quality comparisons, benchmarks, and more.
- Read the [SDXL in 4 steps with Latent Consistency LoRAs](https://huggingface.co/blog/lcm_lora) blog post to learn more about SDXL LCM-LoRA's for super fast inference, quality comparisons, benchmarks, and more.

View File

@@ -77,7 +77,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -113,50 +113,36 @@ The dataset preprocessing code and training loop are found in the [`main()`](htt
As with the script parameters, a walkthrough of the training script is provided in the [Text-to-image](text2image#training-script) training guide. Instead, this guide takes a look at the LoRA relevant parts of the script.
<hfoptions id="lora">
<hfoption id="UNet">
Diffusers uses [`~peft.LoraConfig`] from the [PEFT](https://hf.co/docs/peft) library to set up the parameters of the LoRA adapter such as the rank, alpha, and which modules to insert the LoRA weights into. The adapter is added to the UNet, and only the LoRA layers are filtered for optimization in `lora_layers`.
The script begins by adding the [new LoRA weights](https://github.com/huggingface/diffusers/blob/dd9a5caf61f04d11c0fa9f3947b69ab0010c9a0f/examples/text_to_image/train_text_to_image_lora.py#L447) to the attention layers. This involves correctly configuring the weight size for each block in the UNet. You'll see the `rank` parameter is used to create the [`~models.attention_processor.LoRAAttnProcessor`]:
```py
unet_lora_config = LoraConfig(
r=args.rank,
lora_alpha=args.rank,
init_lora_weights="gaussian",
target_modules=["to_k", "to_q", "to_v", "to_out.0"],
)
lora_attn_procs = {}
for name in unet.attn_processors.keys():
cross_attention_dim = None if name.endswith("attn1.processor") else unet.config.cross_attention_dim
if name.startswith("mid_block"):
hidden_size = unet.config.block_out_channels[-1]
elif name.startswith("up_blocks"):
block_id = int(name[len("up_blocks.")])
hidden_size = list(reversed(unet.config.block_out_channels))[block_id]
elif name.startswith("down_blocks"):
block_id = int(name[len("down_blocks.")])
hidden_size = unet.config.block_out_channels[block_id]
unet.add_adapter(unet_lora_config)
lora_layers = filter(lambda p: p.requires_grad, unet.parameters())
lora_attn_procs[name] = LoRAAttnProcessor(
hidden_size=hidden_size,
cross_attention_dim=cross_attention_dim,
rank=args.rank,
)
unet.set_attn_processor(lora_attn_procs)
lora_layers = AttnProcsLayers(unet.attn_processors)
```
</hfoption>
<hfoption id="text encoder">
Diffusers also supports finetuning the text encoder with LoRA from the [PEFT](https://hf.co/docs/peft) library when necessary such as finetuning Stable Diffusion XL (SDXL). The [`~peft.LoraConfig`] is used to configure the parameters of the LoRA adapter which are then added to the text encoder, and only the LoRA layers are filtered for training.
```py
text_lora_config = LoraConfig(
r=args.rank,
lora_alpha=args.rank,
init_lora_weights="gaussian",
target_modules=["q_proj", "k_proj", "v_proj", "out_proj"],
)
text_encoder_one.add_adapter(text_lora_config)
text_encoder_two.add_adapter(text_lora_config)
text_lora_parameters_one = list(filter(lambda p: p.requires_grad, text_encoder_one.parameters()))
text_lora_parameters_two = list(filter(lambda p: p.requires_grad, text_encoder_two.parameters()))
```
</hfoption>
</hfoptions>
The [optimizer](https://github.com/huggingface/diffusers/blob/e4b8f173b97731686e290b2eb98e7f5df2b1b322/examples/text_to_image/train_text_to_image_lora.py#L529) is initialized with the `lora_layers` because these are the only weights that'll be optimized:
The [optimizer](https://github.com/huggingface/diffusers/blob/dd9a5caf61f04d11c0fa9f3947b69ab0010c9a0f/examples/text_to_image/train_text_to_image_lora.py#L519) is initialized with the `lora_layers` because these are the only weights that'll be optimized:
```py
optimizer = optimizer_cls(
lora_layers,
lora_layers.parameters(),
lr=args.learning_rate,
betas=(args.adam_beta1, args.adam_beta2),
weight_decay=args.adam_weight_decay,
@@ -170,7 +156,7 @@ Aside from setting up the LoRA layers, the training script is more or less the s
Once you've made all your changes or you're okay with the default configuration, you're ready to launch the training script! 🚀
Let's train on the [Pokémon BLIP captions](https://huggingface.co/datasets/lambdalabs/pokemon-blip-captions) dataset to generate our own Pokémon. Set the environment variables `MODEL_NAME` and `DATASET_NAME` to the model and dataset respectively. You should also specify where to save the model in `OUTPUT_DIR`, and the name of the model to save to on the Hub with `HUB_MODEL_ID`. The script creates and saves the following files to your repository:
Let's train on the [Pokémon BLIP captions](https://huggingface.co/datasets/lambdalabs/pokemon-blip-captions) dataset to generate our yown Pokémon. Set the environment variables `MODEL_NAME` and `DATASET_NAME` to the model and dataset respectively. You should also specify where to save the model in `OUTPUT_DIR`, and the name of the model to save to on the Hub with `HUB_MODEL_ID`. The script creates and saves the following files to your repository:
- saved model checkpoints
- `pytorch_lora_weights.safetensors` (the trained LoRA weights)

View File

@@ -59,7 +59,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()

View File

@@ -53,7 +53,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()

View File

@@ -69,7 +69,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()

View File

@@ -67,7 +67,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()

View File

@@ -51,7 +51,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()

View File

@@ -53,7 +53,7 @@ accelerate config default
Or if your environment doesn't support an interactive shell, like a notebook, you can use:
```py
```bash
from accelerate.utils import write_basic_config
write_basic_config()
@@ -173,7 +173,7 @@ pipeline = AutoPipelineForText2Image.from_pretrained("path/to/saved/model", torc
caption = "A cute bird pokemon holding a shield"
images = pipeline(
caption,
caption,
width=1024,
height=1536,
prior_timesteps=DEFAULT_STAGE_C_TIMESTEPS,

View File

@@ -12,74 +12,75 @@ specific language governing permissions and limitations under the License.
# AutoPipeline
Diffusers provides many pipelines for basic tasks like generating images, videos, audio, and inpainting. On top of these, there are specialized pipelines for adapters and features like upscaling, super-resolution, and more. Different pipeline classes can even use the same checkpoint because they share the same pretrained model! With so many different pipelines, it can be overwhelming to know which pipeline class to use.
🤗 Diffusers is able to complete many different tasks, and you can often reuse the same pretrained weights for multiple tasks such as text-to-image, image-to-image, and inpainting. If you're new to the library and diffusion models though, it may be difficult to know which pipeline to use for a task. For example, if you're using the [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5) checkpoint for text-to-image, you might not know that you could also use it for image-to-image and inpainting by loading the checkpoint with the [`StableDiffusionImg2ImgPipeline`] and [`StableDiffusionInpaintPipeline`] classes respectively.
The [AutoPipeline](../api/pipelines/auto_pipeline) class is designed to simplify the variety of pipelines in Diffusers. It is a generic *task-first* pipeline that lets you focus on a task ([`AutoPipelineForText2Image`], [`AutoPipelineForImage2Image`], and [`AutoPipelineForInpainting`]) without needing to know the specific pipeline class. The [AutoPipeline](../api/pipelines/auto_pipeline) automatically detects the correct pipeline class to use.
The `AutoPipeline` class is designed to simplify the variety of pipelines in 🤗 Diffusers. It is a generic, *task-first* pipeline that lets you focus on the task. The `AutoPipeline` automatically detects the correct pipeline class to use, which makes it easier to load a checkpoint for a task without knowing the specific pipeline class name.
For example, let's use the [dreamlike-art/dreamlike-photoreal-2.0](https://hf.co/dreamlike-art/dreamlike-photoreal-2.0) checkpoint.
<Tip>
Under the hood, [AutoPipeline](../api/pipelines/auto_pipeline):
Take a look at the [AutoPipeline](../api/pipelines/auto_pipeline) reference to see which tasks are supported. Currently, it supports text-to-image, image-to-image, and inpainting.
1. Detects a `"stable-diffusion"` class from the [model_index.json](https://hf.co/dreamlike-art/dreamlike-photoreal-2.0/blob/main/model_index.json) file.
2. Depending on the task you're interested in, it loads the [`StableDiffusionPipeline`], [`StableDiffusionImg2ImgPipeline`], or [`StableDiffusionInpaintPipeline`]. Any parameter (`strength`, `num_inference_steps`, etc.) you would pass to these specific pipelines can also be passed to the [AutoPipeline](../api/pipelines/auto_pipeline).
</Tip>
<hfoptions id="autopipeline">
<hfoption id="text-to-image">
This tutorial shows you how to use an `AutoPipeline` to automatically infer the pipeline class to load for a specific task, given the pretrained weights.
## Choose an AutoPipeline for your task
Start by picking a checkpoint. For example, if you're interested in text-to-image with the [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5) checkpoint, use [`AutoPipelineForText2Image`]:
```py
from diffusers import AutoPipelineForText2Image
import torch
pipe_txt2img = AutoPipelineForText2Image.from_pretrained(
"dreamlike-art/dreamlike-photoreal-2.0", torch_dtype=torch.float16, use_safetensors=True
pipeline = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
prompt = "peasant and dragon combat, wood cutting style, viking era, bevel with rune"
prompt = "cinematic photo of Godzilla eating sushi with a cat in a izakaya, 35mm photograph, film, professional, 4k, highly detailed"
generator = torch.Generator(device="cpu").manual_seed(37)
image = pipe_txt2img(prompt, generator=generator).images[0]
image = pipeline(prompt, num_inference_steps=25).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-text2img.png"/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-text2img.png" alt="generated image of peasant fighting dragon in wood cutting style"/>
</div>
</hfoption>
<hfoption id="image-to-image">
Under the hood, [`AutoPipelineForText2Image`]:
1. automatically detects a `"stable-diffusion"` class from the [`model_index.json`](https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/model_index.json) file
2. loads the corresponding text-to-image [`StableDiffusionPipeline`] based on the `"stable-diffusion"` class name
Likewise, for image-to-image, [`AutoPipelineForImage2Image`] detects a `"stable-diffusion"` checkpoint from the `model_index.json` file and it'll load the corresponding [`StableDiffusionImg2ImgPipeline`] behind the scenes. You can also pass any additional arguments specific to the pipeline class such as `strength`, which determines the amount of noise or variation added to an input image:
```py
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import load_image
import torch
import requests
from PIL import Image
from io import BytesIO
pipe_img2img = AutoPipelineForImage2Image.from_pretrained(
"dreamlike-art/dreamlike-photoreal-2.0", torch_dtype=torch.float16, use_safetensors=True
pipeline = AutoPipelineForImage2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
).to("cuda")
prompt = "a portrait of a dog wearing a pearl earring"
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-text2img.png")
url = "https://upload.wikimedia.org/wikipedia/commons/thumb/0/0f/1665_Girl_with_a_Pearl_Earring.jpg/800px-1665_Girl_with_a_Pearl_Earring.jpg"
prompt = "cinematic photo of Godzilla eating burgers with a cat in a fast food restaurant, 35mm photograph, film, professional, 4k, highly detailed"
generator = torch.Generator(device="cpu").manual_seed(53)
image = pipe_img2img(prompt, image=init_image, generator=generator).images[0]
response = requests.get(url)
image = Image.open(BytesIO(response.content)).convert("RGB")
image.thumbnail((768, 768))
image = pipeline(prompt, image, num_inference_steps=200, strength=0.75, guidance_scale=10.5).images[0]
image
```
Notice how the [dreamlike-art/dreamlike-photoreal-2.0](https://hf.co/dreamlike-art/dreamlike-photoreal-2.0) checkpoint is used for both text-to-image and image-to-image tasks? To save memory and avoid loading the checkpoint twice, use the [`~DiffusionPipeline.from_pipe`] method.
```py
pipe_img2img = AutoPipelineForImage2Image.from_pipe(pipe_txt2img).to("cuda")
image = pipeline(prompt, image=init_image, generator=generator).images[0]
image
```
You can learn more about the [`~DiffusionPipeline.from_pipe`] method in the [Reuse a pipeline](../using-diffusers/loading#reuse-a-pipeline) guide.
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-img2img.png"/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-img2img.png" alt="generated image of a vermeer portrait of a dog wearing a pearl earring"/>
</div>
</hfoption>
<hfoption id="inpainting">
And if you want to do inpainting, then [`AutoPipelineForInpainting`] loads the underlying [`StableDiffusionInpaintPipeline`] class in the same way:
```py
from diffusers import AutoPipelineForInpainting
@@ -90,27 +91,22 @@ pipeline = AutoPipelineForInpainting.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-img2img.png")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-mask.png")
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
prompt = "cinematic photo of a owl, 35mm photograph, film, professional, 4k, highly detailed"
generator = torch.Generator(device="cpu").manual_seed(38)
image = pipeline(prompt, image=init_image, mask_image=mask_image, generator=generator, strength=0.4).images[0]
init_image = load_image(img_url).convert("RGB")
mask_image = load_image(mask_url).convert("RGB")
prompt = "A majestic tiger sitting on a bench"
image = pipeline(prompt, image=init_image, mask_image=mask_image, num_inference_steps=50, strength=0.80).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-inpaint.png"/>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/autopipeline-inpaint.png" alt="generated image of a tiger sitting on a bench"/>
</div>
</hfoption>
</hfoptions>
## Unsupported checkpoints
The [AutoPipeline](../api/pipelines/auto_pipeline) supports [Stable Diffusion](../api/pipelines/stable_diffusion/overview), [Stable Diffusion XL](../api/pipelines/stable_diffusion/stable_diffusion_xl), [ControlNet](../api/pipelines/controlnet), [Kandinsky 2.1](../api/pipelines/kandinsky.md), [Kandinsky 2.2](../api/pipelines/kandinsky_v22), and [DeepFloyd IF](../api/pipelines/deepfloyd_if) checkpoints.
If you try to load an unsupported checkpoint, you'll get an error.
If you try to load an unsupported checkpoint, it'll throw an error:
```py
from diffusers import AutoPipelineForImage2Image
@@ -121,3 +117,54 @@ pipeline = AutoPipelineForImage2Image.from_pretrained(
)
"ValueError: AutoPipeline can't find a pipeline linked to ShapEImg2ImgPipeline for None"
```
## Use multiple pipelines
For some workflows or if you're loading many pipelines, it is more memory-efficient to reuse the same components from a checkpoint instead of reloading them which would unnecessarily consume additional memory. For example, if you're using a checkpoint for text-to-image and you want to use it again for image-to-image, use the [`~AutoPipelineForImage2Image.from_pipe`] method. This method creates a new pipeline from the components of a previously loaded pipeline at no additional memory cost.
The [`~AutoPipelineForImage2Image.from_pipe`] method detects the original pipeline class and maps it to the new pipeline class corresponding to the task you want to do. For example, if you load a `"stable-diffusion"` class pipeline for text-to-image:
```py
from diffusers import AutoPipelineForText2Image, AutoPipelineForImage2Image
import torch
pipeline_text2img = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
)
print(type(pipeline_text2img))
"<class 'diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline'>"
```
Then [`~AutoPipelineForImage2Image.from_pipe`] maps the original `"stable-diffusion"` pipeline class to [`StableDiffusionImg2ImgPipeline`]:
```py
pipeline_img2img = AutoPipelineForImage2Image.from_pipe(pipeline_text2img)
print(type(pipeline_img2img))
"<class 'diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_img2img.StableDiffusionImg2ImgPipeline'>"
```
If you passed an optional argument - like disabling the safety checker - to the original pipeline, this argument is also passed on to the new pipeline:
```py
from diffusers import AutoPipelineForText2Image, AutoPipelineForImage2Image
import torch
pipeline_text2img = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
requires_safety_checker=False,
).to("cuda")
pipeline_img2img = AutoPipelineForImage2Image.from_pipe(pipeline_text2img)
print(pipeline_img2img.config.requires_safety_checker)
"False"
```
You can overwrite any of the arguments and even configuration from the original pipeline if you want to change the behavior of the new pipeline. For example, to turn the safety checker back on and add the `strength` argument:
```py
pipeline_img2img = AutoPipelineForImage2Image.from_pipe(pipeline_text2img, requires_safety_checker=True, strength=0.3)
print(pipeline_img2img.config.requires_safety_checker)
"True"
```

View File

@@ -14,17 +14,19 @@ specific language governing permissions and limitations under the License.
# Load LoRAs for inference
There are many adapter types (with [LoRAs](https://huggingface.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora) being the most popular) trained in different styles to achieve different effects. You can even combine multiple adapters to create new and unique images.
There are many adapters (with LoRAs being the most common type) trained in different styles to achieve different effects. You can even combine multiple adapters to create new and unique images. With the 🤗 [PEFT](https://huggingface.co/docs/peft/index) integration in 🤗 Diffusers, it is really easy to load and manage adapters for inference. In this guide, you'll learn how to use different adapters with [Stable Diffusion XL (SDXL)](../api/pipelines/stable_diffusion/stable_diffusion_xl) for inference.
In this tutorial, you'll learn how to easily load and manage adapters for inference with the 🤗 [PEFT](https://huggingface.co/docs/peft/index) integration in 🤗 Diffusers. You'll use LoRA as the main adapter technique, so you'll see the terms LoRA and adapter used interchangeably.
Throughout this guide, you'll use LoRA as the main adapter technique, so we'll use the terms LoRA and adapter interchangeably. You should have some familiarity with LoRA, and if you don't, we welcome you to check out the [LoRA guide](https://huggingface.co/docs/peft/conceptual_guides/lora).
Let's first install all the required libraries.
```bash
!pip install -q transformers accelerate peft diffusers
!pip install -q transformers accelerate
!pip install peft
!pip install diffusers
```
Now, load a pipeline with a [Stable Diffusion XL (SDXL)](../api/pipelines/stable_diffusion/stable_diffusion_xl) checkpoint:
Now, let's load a pipeline with a SDXL checkpoint:
```python
from diffusers import DiffusionPipeline
@@ -34,18 +36,21 @@ pipe_id = "stabilityai/stable-diffusion-xl-base-1.0"
pipe = DiffusionPipeline.from_pretrained(pipe_id, torch_dtype=torch.float16).to("cuda")
```
Next, load a [CiroN2022/toy-face](https://huggingface.co/CiroN2022/toy-face) adapter with the [`~diffusers.loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights`] method. With the 🤗 PEFT integration, you can assign a specific `adapter_name` to the checkpoint, which let's you easily switch between different LoRA checkpoints. Let's call this adapter `"toy"`.
Next, load a LoRA checkpoint with the [`~diffusers.loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights`] method.
With the 🤗 PEFT integration, you can assign a specific `adapter_name` to the checkpoint, which let's you easily switch between different LoRA checkpoints. Let's call this adapter `"toy"`.
```python
pipe.load_lora_weights("CiroN2022/toy-face", weight_name="toy_face_sdxl.safetensors", adapter_name="toy")
```
Make sure to include the token `toy_face` in the prompt and then you can perform inference:
And then perform inference:
```python
prompt = "toy_face of a hacker with a hoodie"
lora_scale = 0.9
lora_scale= 0.9
image = pipe(
prompt, num_inference_steps=30, cross_attention_kwargs={"scale": lora_scale}, generator=torch.manual_seed(0)
).images[0]
@@ -54,16 +59,17 @@ image
![toy-face](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_8_1.png)
With the `adapter_name` parameter, it is really easy to use another adapter for inference! Load the [nerijs/pixel-art-xl](https://huggingface.co/nerijs/pixel-art-xl) adapter that has been fine-tuned to generate pixel art images and call it `"pixel"`.
The pipeline automatically sets the first loaded adapter (`"toy"`) as the active adapter, but you can activate the `"pixel"` adapter with the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method:
With the `adapter_name` parameter, it is really easy to use another adapter for inference! Load the [nerijs/pixel-art-xl](https://huggingface.co/nerijs/pixel-art-xl) adapter that has been fine-tuned to generate pixel art images, and let's call it `"pixel"`.
The pipeline automatically sets the first loaded adapter (`"toy"`) as the active adapter. But you can activate the `"pixel"` adapter with the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method as shown below:
```python
pipe.load_lora_weights("nerijs/pixel-art-xl", weight_name="pixel-art-xl.safetensors", adapter_name="pixel")
pipe.set_adapters("pixel")
```
Make sure you include the token `pixel art` in your prompt to generate a pixel art image:
Let's now generate an image with the second adapter and check the result:
```python
prompt = "a hacker with a hoodie, pixel art"
@@ -75,25 +81,29 @@ image
![pixel-art](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_12_1.png)
## Merge adapters
## Combine multiple adapters
You can also merge different adapter checkpoints for inference to blend their styles together.
You can also perform multi-adapter inference where you combine different adapter checkpoints for inference.
Once again, use the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method to activate the `pixel` and `toy` adapters and specify the weights for how they should be merged.
Once again, use the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method to activate two LoRA checkpoints and specify the weight for how the checkpoints should be combined.
```python
pipe.set_adapters(["pixel", "toy"], adapter_weights=[0.5, 1.0])
```
Now that we have set these two adapters, let's generate an image from the combined adapters!
<Tip>
LoRA checkpoints in the diffusion community are almost always obtained with [DreamBooth](https://huggingface.co/docs/diffusers/main/en/training/dreambooth). DreamBooth training often relies on "trigger" words in the input text prompts in order for the generation results to look as expected. When you combine multiple LoRA checkpoints, it's important to ensure the trigger words for the corresponding LoRA checkpoints are present in the input text prompts.
</Tip>
Remember to use the trigger words for [CiroN2022/toy-face](https://hf.co/CiroN2022/toy-face) and [nerijs/pixel-art-xl](https://hf.co/nerijs/pixel-art-xl) (these are found in their repositories) in the prompt to generate an image.
The trigger words for [CiroN2022/toy-face](https://hf.co/CiroN2022/toy-face) and [nerijs/pixel-art-xl](https://hf.co/nerijs/pixel-art-xl) are found in their repositories.
```python
# Notice how the prompt is constructed.
prompt = "toy_face of a hacker with a hoodie, pixel art"
image = pipe(
prompt, num_inference_steps=30, cross_attention_kwargs={"scale": 1.0}, generator=torch.manual_seed(0)
@@ -103,95 +113,43 @@ image
![toy-face-pixel-art](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_16_1.png)
Impressive! As you can see, the model generated an image that mixed the characteristics of both adapters.
Impressive! As you can see, the model was able to generate an image that mixes the characteristics of both adapters.
> [!TIP]
> Through its PEFT integration, Diffusers also offers more efficient merging methods which you can learn about in the [Merge LoRAs](../using-diffusers/merge_loras) guide!
To return to only using one adapter, use the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method to activate the `"toy"` adapter:
If you want to go back to using only one adapter, use the [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`] method to activate the `"toy"` adapter:
```python
# First, set the adapter.
pipe.set_adapters("toy")
# Then, run inference.
prompt = "toy_face of a hacker with a hoodie"
lora_scale = 0.9
lora_scale= 0.9
image = pipe(
prompt, num_inference_steps=30, cross_attention_kwargs={"scale": lora_scale}, generator=torch.manual_seed(0)
).images[0]
image
```
Or to disable all adapters entirely, use the [`~diffusers.loaders.UNet2DConditionLoadersMixin.disable_lora`] method to return the base model.
![toy-face-again](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_18_1.png)
If you want to switch to only the base model, disable all LoRAs with the [`~diffusers.loaders.UNet2DConditionLoadersMixin.disable_lora`] method.
```python
pipe.disable_lora()
prompt = "toy_face of a hacker with a hoodie"
lora_scale= 0.9
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
image
```
![no-lora](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_20_1.png)
### Customize adapters strength
For even more customization, you can control how strongly the adapter affects each part of the pipeline. For this, pass a dictionary with the control strengths (called "scales") to [`~diffusers.loaders.UNet2DConditionLoadersMixin.set_adapters`].
## Monitoring active adapters
For example, here's how you can turn on the adapter for the `down` parts, but turn it off for the `mid` and `up` parts:
```python
pipe.enable_lora() # enable lora again, after we disabled it above
prompt = "toy_face of a hacker with a hoodie, pixel art"
adapter_weight_scales = { "unet": { "down": 1, "mid": 0, "up": 0} }
pipe.set_adapters("pixel", adapter_weight_scales)
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
image
```
![block-lora-text-and-down](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_block_down.png)
Let's see how turning off the `down` part and turning on the `mid` and `up` part respectively changes the image.
```python
adapter_weight_scales = { "unet": { "down": 0, "mid": 1, "up": 0} }
pipe.set_adapters("pixel", adapter_weight_scales)
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
image
```
![block-lora-text-and-mid](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_block_mid.png)
```python
adapter_weight_scales = { "unet": { "down": 0, "mid": 0, "up": 1} }
pipe.set_adapters("pixel", adapter_weight_scales)
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
image
```
![block-lora-text-and-up](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_block_up.png)
Looks cool!
This is a really powerful feature. You can use it to control the adapter strengths down to per-transformer level. And you can even use it for multiple adapters.
```python
adapter_weight_scales_toy = 0.5
adapter_weight_scales_pixel = {
"unet": {
"down": 0.9, # all transformers in the down-part will use scale 0.9
# "mid" # because, in this example, "mid" is not given, all transformers in the mid part will use the default scale 1.0
"up": {
"block_0": 0.6, # all 3 transformers in the 0th block in the up-part will use scale 0.6
"block_1": [0.4, 0.8, 1.0], # the 3 transformers in the 1st block in the up-part will use scales 0.4, 0.8 and 1.0 respectively
}
}
}
pipe.set_adapters(["toy", "pixel"], [adapter_weight_scales_toy, adapter_weight_scales_pixel])
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
image
```
![block-lora-mixed](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/peft_integration/diffusers_peft_lora_inference_block_mixed.png)
## Manage active adapters
You have attached multiple adapters in this tutorial, and if you're feeling a bit lost on what adapters have been attached to the pipeline's components, use the [`~diffusers.loaders.LoraLoaderMixin.get_active_adapters`] method to check the list of active adapters:
You have attached multiple adapters in this tutorial, and if you're feeling a bit lost on what adapters have been attached to the pipeline's components, you can easily check the list of active adapters using the [`~diffusers.loaders.LoraLoaderMixin.get_active_adapters`] method:
```py
active_adapters = pipe.get_active_adapters()
@@ -206,3 +164,74 @@ list_adapters_component_wise = pipe.get_list_adapters()
list_adapters_component_wise
{"text_encoder": ["toy", "pixel"], "unet": ["toy", "pixel"], "text_encoder_2": ["toy", "pixel"]}
```
## Compatibility with `torch.compile`
If you want to compile your model with `torch.compile` make sure to first fuse the LoRA weights into the base model and unload them.
```py
pipe.load_lora_weights("nerijs/pixel-art-xl", weight_name="pixel-art-xl.safetensors", adapter_name="pixel")
pipe.load_lora_weights("CiroN2022/toy-face", weight_name="toy_face_sdxl.safetensors", adapter_name="toy")
pipe.set_adapters(["pixel", "toy"], adapter_weights=[0.5, 1.0])
# Fuses the LoRAs into the Unet
pipe.fuse_lora()
pipe.unload_lora_weights()
pipe = torch.compile(pipe)
prompt = "toy_face of a hacker with a hoodie, pixel art"
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
```
## Fusing adapters into the model
You can use PEFT to easily fuse/unfuse multiple adapters directly into the model weights (both UNet and text encoder) using the [`~diffusers.loaders.LoraLoaderMixin.fuse_lora`] method, which can lead to a speed-up in inference and lower VRAM usage.
```py
pipe.load_lora_weights("nerijs/pixel-art-xl", weight_name="pixel-art-xl.safetensors", adapter_name="pixel")
pipe.load_lora_weights("CiroN2022/toy-face", weight_name="toy_face_sdxl.safetensors", adapter_name="toy")
pipe.set_adapters(["pixel", "toy"], adapter_weights=[0.5, 1.0])
# Fuses the LoRAs into the Unet
pipe.fuse_lora()
prompt = "toy_face of a hacker with a hoodie, pixel art"
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
# Gets the Unet back to the original state
pipe.unfuse_lora()
```
You can also fuse some adapters using `adapter_names` for faster generation:
```py
pipe.load_lora_weights("nerijs/pixel-art-xl", weight_name="pixel-art-xl.safetensors", adapter_name="pixel")
pipe.load_lora_weights("CiroN2022/toy-face", weight_name="toy_face_sdxl.safetensors", adapter_name="toy")
pipe.set_adapters(["pixel"], adapter_weights=[0.5, 1.0])
# Fuses the LoRAs into the Unet
pipe.fuse_lora(adapter_names=["pixel"])
prompt = "a hacker with a hoodie, pixel art"
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
# Gets the Unet back to the original state
pipe.unfuse_lora()
# Fuse all adapters
pipe.fuse_lora(adapter_names=["pixel", "toy"])
prompt = "toy_face of a hacker with a hoodie, pixel art"
image = pipe(prompt, num_inference_steps=30, generator=torch.manual_seed(0)).images[0]
```
## Saving a pipeline after fusing the adapters
To properly save a pipeline after it's been loaded with the adapters, it should be serialized like so:
```python
pipe.fuse_lora(lora_scale=1.0)
pipe.unload_lora_weights()
pipe.save_pretrained("path-to-pipeline")
```

View File

@@ -12,18 +12,13 @@ specific language governing permissions and limitations under the License.
# Pipeline callbacks
The denoising loop of a pipeline can be modified with custom defined functions using the `callback_on_step_end` parameter. The callback function is executed at the end of each step, and modifies the pipeline attributes and variables for the next step. This is really useful for *dynamically* adjusting certain pipeline attributes or modifying tensor variables. This versatility allows for interesting use-cases such as changing the prompt embeddings at each timestep, assigning different weights to the prompt embeddings, and editing the guidance scale. With callbacks, you can implement new features without modifying the underlying code!
The denoising loop of a pipeline can be modified with custom defined functions using the `callback_on_step_end` parameter. This can be really useful for *dynamically* adjusting certain pipeline attributes, or modifying tensor variables. The flexibility of callbacks opens up some interesting use-cases such as changing the prompt embeddings at each timestep, assigning different weights to the prompt embeddings, and editing the guidance scale.
> [!TIP]
> 🤗 Diffusers currently only supports `callback_on_step_end`, but feel free to open a [feature request](https://github.com/huggingface/diffusers/issues/new/choose) if you have a cool use-case and require a callback function with a different execution point!
This guide will show you how to use the `callback_on_step_end` parameter to disable classifier-free guidance (CFG) after 40% of the inference steps to save compute with minimal cost to performance.
This guide will demonstrate how callbacks work by a few features you can implement with them.
The callback function should have the following arguments:
## Dynamic classifier-free guidance
Dynamic classifier-free guidance (CFG) is a feature that allows you to disable CFG after a certain number of inference steps which can help you save compute with minimal cost to performance. The callback function for this should have the following arguments:
* `pipeline` (or the pipeline instance) provides access to important properties such as `num_timesteps` and `guidance_scale`. You can modify these properties by updating the underlying attributes. For this example, you'll disable CFG by setting `pipeline._guidance_scale=0.0`.
* `pipe` (or the pipeline instance) provides access to useful properties such as `num_timesteps` and `guidance_scale`. You can modify these properties by updating the underlying attributes. For this example, you'll disable CFG by setting `pipe._guidance_scale=0.0`.
* `step_index` and `timestep` tell you where you are in the denoising loop. Use `step_index` to turn off CFG after reaching 40% of `num_timesteps`.
* `callback_kwargs` is a dict that contains tensor variables you can modify during the denoising loop. It only includes variables specified in the `callback_on_step_end_tensor_inputs` argument, which is passed to the pipeline's `__call__` method. Different pipelines may use different sets of variables, so please check a pipeline's `_callback_tensor_inputs` attribute for the list of variables you can modify. Some common variables include `latents` and `prompt_embeds`. For this function, change the batch size of `prompt_embeds` after setting `guidance_scale=0.0` in order for it to work properly.
@@ -32,13 +27,13 @@ Your callback function should look something like this:
```python
def callback_dynamic_cfg(pipe, step_index, timestep, callback_kwargs):
# adjust the batch_size of prompt_embeds according to guidance_scale
if step_index == int(pipeline.num_timesteps * 0.4):
if step_index == int(pipe.num_timesteps * 0.4):
prompt_embeds = callback_kwargs["prompt_embeds"]
prompt_embeds = prompt_embeds.chunk(2)[-1]
# update guidance_scale and prompt_embeds
pipeline._guidance_scale = 0.0
callback_kwargs["prompt_embeds"] = prompt_embeds
# update guidance_scale and prompt_embeds
pipe._guidance_scale = 0.0
callback_kwargs["prompt_embeds"] = prompt_embeds
return callback_kwargs
```
@@ -48,134 +43,58 @@ Now, you can pass the callback function to the `callback_on_step_end` parameter
import torch
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipeline = pipeline.to("cuda")
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipe = pipe.to("cuda")
prompt = "a photo of an astronaut riding a horse on mars"
generator = torch.Generator(device="cuda").manual_seed(1)
out = pipeline(
prompt,
generator=generator,
callback_on_step_end=callback_dynamic_cfg,
callback_on_step_end_tensor_inputs=['prompt_embeds']
)
out = pipe(prompt, generator=generator, callback_on_step_end=callback_dynamic_cfg, callback_on_step_end_tensor_inputs=['prompt_embeds'])
out.images[0].save("out_custom_cfg.png")
```
The callback function is executed at the end of each denoising step, and modifies the pipeline attributes and tensor variables for the next denoising step.
With callbacks, you can implement features such as dynamic CFG without having to modify the underlying code at all!
<Tip>
🤗 Diffusers currently only supports `callback_on_step_end`, but feel free to open a [feature request](https://github.com/huggingface/diffusers/issues/new/choose) if you have a cool use-case and require a callback function with a different execution point!
</Tip>
## Interrupt the diffusion process
> [!TIP]
> The interruption callback is supported for text-to-image, image-to-image, and inpainting for the [StableDiffusionPipeline](../api/pipelines/stable_diffusion/overview) and [StableDiffusionXLPipeline](../api/pipelines/stable_diffusion/stable_diffusion_xl).
Interrupting the diffusion process is particularly useful when building UIs that work with Diffusers because it allows users to stop the generation process if they're unhappy with the intermediate results. You can incorporate this into your pipeline with a callback.
Stopping the diffusion process early is useful when building UIs that work with Diffusers because it allows users to stop the generation process if they're unhappy with the intermediate results. You can incorporate this into your pipeline with a callback.
<Tip>
This callback function should take the following arguments: `pipeline`, `i`, `t`, and `callback_kwargs` (this must be returned). Set the pipeline's `_interrupt` attribute to `True` to stop the diffusion process after a certain number of steps. You are also free to implement your own custom stopping logic inside the callback.
The interruption callback is supported for text-to-image, image-to-image, and inpainting for the [StableDiffusionPipeline](../api/pipelines/stable_diffusion/overview) and [StableDiffusionXLPipeline](../api/pipelines/stable_diffusion/stable_diffusion_xl).
</Tip>
This callback function should take the following arguments: `pipe`, `i`, `t`, and `callback_kwargs` (this must be returned). Set the pipeline's `_interrupt` attribute to `True` to stop the diffusion process after a certain number of steps. You are also free to implement your own custom stopping logic inside the callback.
In this example, the diffusion process is stopped after 10 steps even though `num_inference_steps` is set to 50.
```python
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
pipeline.enable_model_cpu_offload()
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5")
pipe.enable_model_cpu_offload()
num_inference_steps = 50
def interrupt_callback(pipeline, i, t, callback_kwargs):
def interrupt_callback(pipe, i, t, callback_kwargs):
stop_idx = 10
if i == stop_idx:
pipeline._interrupt = True
pipe._interrupt = True
return callback_kwargs
pipeline(
pipe(
"A photo of a cat",
num_inference_steps=num_inference_steps,
callback_on_step_end=interrupt_callback,
)
```
## Display image after each generation step
> [!TIP]
> This tip was contributed by [asomoza](https://github.com/asomoza).
Display an image after each generation step by accessing and converting the latents after each step into an image. The latent space is compressed to 128x128, so the images are also 128x128 which is useful for a quick preview.
1. Use the function below to convert the SDXL latents (4 channels) to RGB tensors (3 channels) as explained in the [Explaining the SDXL latent space](https://huggingface.co/blog/TimothyAlexisVass/explaining-the-sdxl-latent-space) blog post.
```py
def latents_to_rgb(latents):
weights = (
(60, -60, 25, -70),
(60, -5, 15, -50),
(60, 10, -5, -35)
)
weights_tensor = torch.t(torch.tensor(weights, dtype=latents.dtype).to(latents.device))
biases_tensor = torch.tensor((150, 140, 130), dtype=latents.dtype).to(latents.device)
rgb_tensor = torch.einsum("...lxy,lr -> ...rxy", latents, weights_tensor) + biases_tensor.unsqueeze(-1).unsqueeze(-1)
image_array = rgb_tensor.clamp(0, 255)[0].byte().cpu().numpy()
image_array = image_array.transpose(1, 2, 0)
return Image.fromarray(image_array)
```
2. Create a function to decode and save the latents into an image.
```py
def decode_tensors(pipe, step, timestep, callback_kwargs):
latents = callback_kwargs["latents"]
image = latents_to_rgb(latents)
image.save(f"{step}.png")
return callback_kwargs
```
3. Pass the `decode_tensors` function to the `callback_on_step_end` parameter to decode the tensors after each step. You also need to specify what you want to modify in the `callback_on_step_end_tensor_inputs` parameter, which in this case are the latents.
```py
from diffusers import AutoPipelineForText2Image
import torch
from PIL import Image
pipeline = AutoPipelineForText2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
variant="fp16",
use_safetensors=True
).to("cuda")
image = pipeline(
prompt="A croissant shaped like a cute bear.",
negative_prompt="Deformed, ugly, bad anatomy",
callback_on_step_end=decode_tensors,
callback_on_step_end_tensor_inputs=["latents"],
).images[0]
```
<div class="flex gap-4 justify-center">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/tips_step_0.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">step 0</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/tips_step_19.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">step 19
</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/tips_step_29.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">step 29</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/tips_step_39.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">step 39</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/tips_step_49.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">step 49</figcaption>
</div>
</div>

View File

@@ -0,0 +1,184 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Contribute a community pipeline
<Tip>
💡 Take a look at GitHub Issue [#841](https://github.com/huggingface/diffusers/issues/841) for more context about why we're adding community pipelines to help everyone easily share their work without being slowed down.
</Tip>
Community pipelines allow you to add any additional features you'd like on top of the [`DiffusionPipeline`]. The main benefit of building on top of the `DiffusionPipeline` is anyone can load and use your pipeline by only adding one more argument, making it super easy for the community to access.
This guide will show you how to create a community pipeline and explain how they work. To keep things simple, you'll create a "one-step" pipeline where the `UNet` does a single forward pass and calls the scheduler once.
## Initialize the pipeline
You should start by creating a `one_step_unet.py` file for your community pipeline. In this file, create a pipeline class that inherits from the [`DiffusionPipeline`] to be able to load model weights and the scheduler configuration from the Hub. The one-step pipeline needs a `UNet` and a scheduler, so you'll need to add these as arguments to the `__init__` function:
```python
from diffusers import DiffusionPipeline
import torch
class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
def __init__(self, unet, scheduler):
super().__init__()
```
To ensure your pipeline and its components (`unet` and `scheduler`) can be saved with [`~DiffusionPipeline.save_pretrained`], add them to the `register_modules` function:
```diff
from diffusers import DiffusionPipeline
import torch
class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
def __init__(self, unet, scheduler):
super().__init__()
+ self.register_modules(unet=unet, scheduler=scheduler)
```
Cool, the `__init__` step is done and you can move to the forward pass now! 🔥
## Define the forward pass
In the forward pass, which we recommend defining as `__call__`, you have complete creative freedom to add whatever feature you'd like. For our amazing one-step pipeline, create a random image and only call the `unet` and `scheduler` once by setting `timestep=1`:
```diff
from diffusers import DiffusionPipeline
import torch
class UnetSchedulerOneForwardPipeline(DiffusionPipeline):
def __init__(self, unet, scheduler):
super().__init__()
self.register_modules(unet=unet, scheduler=scheduler)
+ def __call__(self):
+ image = torch.randn(
+ (1, self.unet.config.in_channels, self.unet.config.sample_size, self.unet.config.sample_size),
+ )
+ timestep = 1
+ model_output = self.unet(image, timestep).sample
+ scheduler_output = self.scheduler.step(model_output, timestep, image).prev_sample
+ return scheduler_output
```
That's it! 🚀 You can now run this pipeline by passing a `unet` and `scheduler` to it:
```python
from diffusers import DDPMScheduler, UNet2DModel
scheduler = DDPMScheduler()
unet = UNet2DModel()
pipeline = UnetSchedulerOneForwardPipeline(unet=unet, scheduler=scheduler)
output = pipeline()
```
But what's even better is you can load pre-existing weights into the pipeline if the pipeline structure is identical. For example, you can load the [`google/ddpm-cifar10-32`](https://huggingface.co/google/ddpm-cifar10-32) weights into the one-step pipeline:
```python
pipeline = UnetSchedulerOneForwardPipeline.from_pretrained("google/ddpm-cifar10-32", use_safetensors=True)
output = pipeline()
```
## Share your pipeline
Open a Pull Request on the 🧨 Diffusers [repository](https://github.com/huggingface/diffusers) to add your awesome pipeline in `one_step_unet.py` to the [examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community) subfolder.
Once it is merged, anyone with `diffusers >= 0.4.0` installed can use this pipeline magically 🪄 by specifying it in the `custom_pipeline` argument:
```python
from diffusers import DiffusionPipeline
pipe = DiffusionPipeline.from_pretrained(
"google/ddpm-cifar10-32", custom_pipeline="one_step_unet", use_safetensors=True
)
pipe()
```
Another way to share your community pipeline is to upload the `one_step_unet.py` file directly to your preferred [model repository](https://huggingface.co/docs/hub/models-uploading) on the Hub. Instead of specifying the `one_step_unet.py` file, pass the model repository id to the `custom_pipeline` argument:
```python
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"google/ddpm-cifar10-32", custom_pipeline="stevhliu/one_step_unet", use_safetensors=True
)
```
Take a look at the following table to compare the two sharing workflows to help you decide the best option for you:
| | GitHub community pipeline | HF Hub community pipeline |
|----------------|------------------------------------------------------------------------------------------------------------------|-------------------------------------------------------------------------------------------|
| usage | same | same |
| review process | open a Pull Request on GitHub and undergo a review process from the Diffusers team before merging; may be slower | upload directly to a Hub repository without any review; this is the fastest workflow |
| visibility | included in the official Diffusers repository and documentation | included on your HF Hub profile and relies on your own usage/promotion to gain visibility |
<Tip>
💡 You can use whatever package you want in your community pipeline file - as long as the user has it installed, everything will work fine. Make sure you have one and only one pipeline class that inherits from `DiffusionPipeline` because this is automatically detected.
</Tip>
## How do community pipelines work?
A community pipeline is a class that inherits from [`DiffusionPipeline`] which means:
- It can be loaded with the [`custom_pipeline`] argument.
- The model weights and scheduler configuration are loaded from [`pretrained_model_name_or_path`].
- The code that implements a feature in the community pipeline is defined in a `pipeline.py` file.
Sometimes you can't load all the pipeline components weights from an official repository. In this case, the other components should be passed directly to the pipeline:
```python
from diffusers import DiffusionPipeline
from transformers import CLIPImageProcessor, CLIPModel
model_id = "CompVis/stable-diffusion-v1-4"
clip_model_id = "laion/CLIP-ViT-B-32-laion2B-s34B-b79K"
feature_extractor = CLIPImageProcessor.from_pretrained(clip_model_id)
clip_model = CLIPModel.from_pretrained(clip_model_id, torch_dtype=torch.float16)
pipeline = DiffusionPipeline.from_pretrained(
model_id,
custom_pipeline="clip_guided_stable_diffusion",
clip_model=clip_model,
feature_extractor=feature_extractor,
scheduler=scheduler,
torch_dtype=torch.float16,
use_safetensors=True,
)
```
The magic behind community pipelines is contained in the following code. It allows the community pipeline to be loaded from GitHub or the Hub, and it'll be available to all 🧨 Diffusers packages.
```python
# 2. Load the pipeline class, if using custom module then load it from the Hub
# if we load from explicit class, let's use it
if custom_pipeline is not None:
pipeline_class = get_class_from_dynamic_module(
custom_pipeline, module_file=CUSTOM_PIPELINE_FILE_NAME, cache_dir=custom_pipeline
)
elif cls != DiffusionPipeline:
pipeline_class = cls
else:
diffusers_module = importlib.import_module(cls.__module__.split(".")[0])
pipeline_class = getattr(diffusers_module, config_dict["_class_name"])
```

View File

@@ -0,0 +1,58 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Control image brightness
The Stable Diffusion pipeline is mediocre at generating images that are either very bright or dark as explained in the [Common Diffusion Noise Schedules and Sample Steps are Flawed](https://huggingface.co/papers/2305.08891) paper. The solutions proposed in the paper are currently implemented in the [`DDIMScheduler`] which you can use to improve the lighting in your images.
<Tip>
💡 Take a look at the paper linked above for more details about the proposed solutions!
</Tip>
One of the solutions is to train a model with *v prediction* and *v loss*. Add the following flag to the [`train_text_to_image.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image.py) or [`train_text_to_image_lora.py`](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image_lora.py) scripts to enable `v_prediction`:
```bash
--prediction_type="v_prediction"
```
For example, let's use the [`ptx0/pseudo-journey-v2`](https://huggingface.co/ptx0/pseudo-journey-v2) checkpoint which has been finetuned with `v_prediction`.
Next, configure the following parameters in the [`DDIMScheduler`]:
1. `rescale_betas_zero_snr=True`, rescales the noise schedule to zero terminal signal-to-noise ratio (SNR)
2. `timestep_spacing="trailing"`, starts sampling from the last timestep
```py
from diffusers import DiffusionPipeline, DDIMScheduler
pipeline = DiffusionPipeline.from_pretrained("ptx0/pseudo-journey-v2", use_safetensors=True)
# switch the scheduler in the pipeline to use the DDIMScheduler
pipeline.scheduler = DDIMScheduler.from_config(
pipeline.scheduler.config, rescale_betas_zero_snr=True, timestep_spacing="trailing"
)
pipeline.to("cuda")
```
Finally, in your call to the pipeline, set `guidance_rescale` to prevent overexposure:
```py
prompt = "A lion in galaxies, spirals, nebulae, stars, smoke, iridescent, intricate detail, octane render, 8k"
image = pipeline(prompt, guidance_rescale=0.7).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/zero_snr.png"/>
</div>

View File

@@ -429,27 +429,6 @@ image = pipe(
make_image_grid([original_image, canny_image, image], rows=1, cols=3)
```
<Tip>
You can use a refiner model with `StableDiffusionXLControlNetPipeline` to improve image quality, just like you can with a regular `StableDiffusionXLPipeline`.
See the [Refine image quality](./sdxl#refine-image-quality) section to learn how to use the refiner model.
Make sure to use `StableDiffusionXLControlNetPipeline` and pass `image` and `controlnet_conditioning_scale`.
```py
base = StableDiffusionXLControlNetPipeline(...)
image = base(
prompt=prompt,
controlnet_conditioning_scale=0.5,
image=canny_image,
num_inference_steps=40,
denoising_end=0.8,
output_type="latent",
).images
# rest exactly as with StableDiffusionXLPipeline
```
</Tip>
## MultiControlNet
<Tip>

View File

@@ -0,0 +1,119 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Community pipelines
[[open-in-colab]]
<Tip>
For more context about the design choices behind community pipelines, please have a look at [this issue](https://github.com/huggingface/diffusers/issues/841).
</Tip>
Community pipelines allow you to get creative and build your own unique pipelines to share with the community. You can find all community pipelines in the [diffusers/examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community) folder along with inference and training examples for how to use them. This guide showcases some of the community pipelines and hopefully it'll inspire you to create your own (feel free to open a PR with your own pipeline and we will merge it!).
To load a community pipeline, use the `custom_pipeline` argument in [`DiffusionPipeline`] to specify one of the files in [diffusers/examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community):
```py
from diffusers import DiffusionPipeline
pipe = DiffusionPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4", custom_pipeline="filename_in_the_community_folder", use_safetensors=True
)
```
If a community pipeline doesn't work as expected, please open a GitHub issue and mention the author.
You can learn more about community pipelines in the how to [load community pipelines](custom_pipeline_overview) and how to [contribute a community pipeline](contribute_pipeline) guides.
## Multilingual Stable Diffusion
The multilingual Stable Diffusion pipeline uses a pretrained [XLM-RoBERTa](https://huggingface.co/papluca/xlm-roberta-base-language-detection) to identify a language and the [mBART-large-50](https://huggingface.co/facebook/mbart-large-50-many-to-one-mmt) model to handle the translation. This allows you to generate images from text in 20 languages.
```py
import torch
from diffusers import DiffusionPipeline
from diffusers.utils import make_image_grid
from transformers import (
pipeline,
MBart50TokenizerFast,
MBartForConditionalGeneration,
)
device = "cuda" if torch.cuda.is_available() else "cpu"
device_dict = {"cuda": 0, "cpu": -1}
# add language detection pipeline
language_detection_model_ckpt = "papluca/xlm-roberta-base-language-detection"
language_detection_pipeline = pipeline("text-classification",
model=language_detection_model_ckpt,
device=device_dict[device])
# add model for language translation
translation_tokenizer = MBart50TokenizerFast.from_pretrained("facebook/mbart-large-50-many-to-one-mmt")
translation_model = MBartForConditionalGeneration.from_pretrained("facebook/mbart-large-50-many-to-one-mmt").to(device)
diffuser_pipeline = DiffusionPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
custom_pipeline="multilingual_stable_diffusion",
detection_pipeline=language_detection_pipeline,
translation_model=translation_model,
translation_tokenizer=translation_tokenizer,
torch_dtype=torch.float16,
)
diffuser_pipeline.enable_attention_slicing()
diffuser_pipeline = diffuser_pipeline.to(device)
prompt = ["a photograph of an astronaut riding a horse",
"Una casa en la playa",
"Ein Hund, der Orange isst",
"Un restaurant parisien"]
images = diffuser_pipeline(prompt).images
make_image_grid(images, rows=2, cols=2)
```
<div class="flex justify-center">
<img src="https://user-images.githubusercontent.com/4313860/198328706-295824a4-9856-4ce5-8e66-278ceb42fd29.png"/>
</div>
## MagicMix
[MagicMix](https://huggingface.co/papers/2210.16056) is a pipeline that can mix an image and text prompt to generate a new image that preserves the image structure. The `mix_factor` determines how much influence the prompt has on the layout generation, `kmin` controls the number of steps during the content generation process, and `kmax` determines how much information is kept in the layout of the original image.
```py
from diffusers import DiffusionPipeline, DDIMScheduler
from diffusers.utils import load_image, make_image_grid
pipeline = DiffusionPipeline.from_pretrained(
"CompVis/stable-diffusion-v1-4",
custom_pipeline="magic_mix",
scheduler=DDIMScheduler.from_pretrained("CompVis/stable-diffusion-v1-4", subfolder="scheduler"),
).to('cuda')
img = load_image("https://user-images.githubusercontent.com/59410571/209578593-141467c7-d831-4792-8b9a-b17dc5e47816.jpg")
mix_img = pipeline(img, prompt="bed", kmin=0.3, kmax=0.5, mix_factor=0.5)
make_image_grid([img, mix_img], rows=1, cols=2)
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://user-images.githubusercontent.com/59410571/209578593-141467c7-d831-4792-8b9a-b17dc5e47816.jpg" />
<figcaption class="mt-2 text-center text-sm text-gray-500">original image</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://user-images.githubusercontent.com/59410571/209578602-70f323fa-05b7-4dd6-b055-e40683e37914.jpg" />
<figcaption class="mt-2 text-center text-sm text-gray-500">image and text prompt mix</figcaption>
</div>
</div>

View File

@@ -16,27 +16,17 @@ specific language governing permissions and limitations under the License.
## Community pipelines
> [!TIP] Take a look at GitHub Issue [#841](https://github.com/huggingface/diffusers/issues/841) for more context about why we're adding community pipelines to help everyone easily share their work without being slowed down.
Community pipelines are any [`DiffusionPipeline`] class that are different from the original implementation as specified in their paper (for example, the [`StableDiffusionControlNetPipeline`] corresponds to the [Text-to-Image Generation with ControlNet Conditioning](https://arxiv.org/abs/2302.05543) paper). They provide additional functionality or extend the original implementation of a pipeline.
Community pipelines are any [`DiffusionPipeline`] class that are different from the original paper implementation (for example, the [`StableDiffusionControlNetPipeline`] corresponds to the [Text-to-Image Generation with ControlNet Conditioning](https://arxiv.org/abs/2302.05543) paper). They provide additional functionality or extend the original implementation of a pipeline.
There are many cool community pipelines like [Speech to Image](https://github.com/huggingface/diffusers/tree/main/examples/community#speech-to-image) or [Composable Stable Diffusion](https://github.com/huggingface/diffusers/tree/main/examples/community#composable-stable-diffusion), and you can find all the official community pipelines [here](https://github.com/huggingface/diffusers/tree/main/examples/community).
There are many cool community pipelines like [Marigold Depth Estimation](https://github.com/huggingface/diffusers/tree/main/examples/community#marigold-depth-estimation) or [InstantID](https://github.com/huggingface/diffusers/tree/main/examples/community#instantid-pipeline), and you can find all the official community pipelines [here](https://github.com/huggingface/diffusers/tree/main/examples/community).
To load any community pipeline on the Hub, pass the repository id of the community pipeline to the `custom_pipeline` argument and the model repository where you'd like to load the pipeline weights and components from. For example, the example below loads a dummy pipeline from [`hf-internal-testing/diffusers-dummy-pipeline`](https://huggingface.co/hf-internal-testing/diffusers-dummy-pipeline/blob/main/pipeline.py) and the pipeline weights and components from [`google/ddpm-cifar10-32`](https://huggingface.co/google/ddpm-cifar10-32):
There are two types of community pipelines, those stored on the Hugging Face Hub and those stored on Diffusers GitHub repository. Hub pipelines are completely customizable (scheduler, models, pipeline code, etc.) while Diffusers GitHub pipelines are only limited to custom pipeline code.
<Tip warning={true}>
| | GitHub community pipeline | HF Hub community pipeline |
|----------------|------------------------------------------------------------------------------------------------------------------|-------------------------------------------------------------------------------------------|
| usage | same | same |
| review process | open a Pull Request on GitHub and undergo a review process from the Diffusers team before merging; may be slower | upload directly to a Hub repository without any review; this is the fastest workflow |
| visibility | included in the official Diffusers repository and documentation | included on your HF Hub profile and relies on your own usage/promotion to gain visibility |
🔒 By loading a community pipeline from the Hugging Face Hub, you are trusting that the code you are loading is safe. Make sure to inspect the code online before loading and running it automatically!
<hfoptions id="community">
<hfoption id="Hub pipelines">
To load a Hugging Face Hub community pipeline, pass the repository id of the community pipeline to the `custom_pipeline` argument and the model repository where you'd like to load the pipeline weights and components from. For example, the example below loads a dummy pipeline from [hf-internal-testing/diffusers-dummy-pipeline](https://huggingface.co/hf-internal-testing/diffusers-dummy-pipeline/blob/main/pipeline.py) and the pipeline weights and components from [google/ddpm-cifar10-32](https://huggingface.co/google/ddpm-cifar10-32):
> [!WARNING]
> By loading a community pipeline from the Hugging Face Hub, you are trusting that the code you are loading is safe. Make sure to inspect the code online before loading and running it automatically!
</Tip>
```py
from diffusers import DiffusionPipeline
@@ -46,10 +36,7 @@ pipeline = DiffusionPipeline.from_pretrained(
)
```
</hfoption>
<hfoption id="GitHub pipelines">
To load a GitHub community pipeline, pass the repository id of the community pipeline to the `custom_pipeline` argument and the model repository where you you'd like to load the pipeline weights and components from. You can also load model components directly. The example below loads the community [CLIP Guided Stable Diffusion](https://github.com/huggingface/diffusers/tree/main/examples/community#clip-guided-stable-diffusion) pipeline and the CLIP model components.
Loading an official community pipeline is similar, but you can mix loading weights from an official repository id and pass pipeline components directly. The example below loads the community [CLIP Guided Stable Diffusion](https://github.com/huggingface/diffusers/tree/main/examples/community#clip-guided-stable-diffusion) pipeline, and you can pass the CLIP model components directly to it:
```py
from diffusers import DiffusionPipeline
@@ -69,12 +56,9 @@ pipeline = DiffusionPipeline.from_pretrained(
)
```
</hfoption>
</hfoptions>
### Load from a local file
Community pipelines can also be loaded from a local file if you pass a file path instead. The path to the passed directory must contain a pipeline.py file that contains the pipeline class.
Community pipelines can also be loaded from a local file if you pass a file path instead. The path to the passed directory must contain a `pipeline.py` file that contains the pipeline class in order to successfully load it.
```py
pipeline = DiffusionPipeline.from_pretrained(
@@ -93,7 +77,7 @@ By default, community pipelines are loaded from the latest stable version of Dif
<hfoptions id="version">
<hfoption id="main">
For example, to load from the main branch:
For example, to load from the `main` branch:
```py
pipeline = DiffusionPipeline.from_pretrained(
@@ -109,7 +93,7 @@ pipeline = DiffusionPipeline.from_pretrained(
</hfoption>
<hfoption id="older version">
For example, to load from a previous version of Diffusers like v0.25.0:
For example, to load from a previous version of Diffusers like `v0.25.0`:
```py
pipeline = DiffusionPipeline.from_pretrained(
@@ -125,140 +109,8 @@ pipeline = DiffusionPipeline.from_pretrained(
</hfoption>
</hfoptions>
### Load with from_pipe
Community pipelines can also be loaded with the [`~DiffusionPipeline.from_pipe`] method which allows you to load and reuse multiple pipelines without any additional memory overhead (learn more in the [Reuse a pipeline](./loading#reuse-a-pipeline) guide). The memory requirement is determined by the largest single pipeline loaded.
For example, let's load a community pipeline that supports [long prompts with weighting](https://github.com/huggingface/diffusers/tree/main/examples/community#long-prompt-weighting-stable-diffusion) from a Stable Diffusion pipeline.
```py
import torch
from diffusers import DiffusionPipeline
pipe_sd = DiffusionPipeline.from_pretrained("emilianJR/CyberRealistic_V3", torch_dtype=torch.float16)
pipe_sd.to("cuda")
# load long prompt weighting pipeline
pipe_lpw = DiffusionPipeline.from_pipe(
pipe_sd,
custom_pipeline="lpw_stable_diffusion",
).to("cuda")
prompt = "cat, hiding in the leaves, ((rain)), zazie rainyday, beautiful eyes, macro shot, colorful details, natural lighting, amazing composition, subsurface scattering, amazing textures, filmic, soft light, ultra-detailed eyes, intricate details, detailed texture, light source contrast, dramatic shadows, cinematic light, depth of field, film grain, noise, dark background, hyperrealistic dslr film still, dim volumetric cinematic lighting"
neg_prompt = "(deformed iris, deformed pupils, semi-realistic, cgi, 3d, render, sketch, cartoon, drawing, anime, mutated hands and fingers:1.4), (deformed, distorted, disfigured:1.3), poorly drawn, bad anatomy, wrong anatomy, extra limb, missing limb, floating limbs, disconnected limbs, mutation, mutated, ugly, disgusting, amputation"
generator = torch.Generator(device="cpu").manual_seed(20)
out_lpw = pipe_lpw(
prompt,
negative_prompt=neg_prompt,
width=512,
height=512,
max_embeddings_multiples=3,
num_inference_steps=50,
generator=generator,
).images[0]
out_lpw
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/from_pipe_lpw.png" />
<figcaption class="mt-2 text-center text-sm text-gray-500">Stable Diffusion with long prompt weighting</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/from_pipe_non_lpw.png" />
<figcaption class="mt-2 text-center text-sm text-gray-500">Stable Diffusion</figcaption>
</div>
</div>
## Example community pipelines
Community pipelines are a really fun and creative way to extend the capabilities of the original pipeline with new and unique features. You can find all community pipelines in the [diffusers/examples/community](https://github.com/huggingface/diffusers/tree/main/examples/community) folder with inference and training examples for how to use them.
This section showcases a couple of the community pipelines and hopefully it'll inspire you to create your own (feel free to open a PR for your community pipeline and ping us for a review)!
> [!TIP]
> The [`~DiffusionPipeline.from_pipe`] method is particularly useful for loading community pipelines because many of them don't have pretrained weights and add a feature on top of an existing pipeline like Stable Diffusion or Stable Diffusion XL. You can learn more about the [`~DiffusionPipeline.from_pipe`] method in the [Load with from_pipe](custom_pipeline_overview#load-with-from_pipe) section.
<hfoptions id="community">
<hfoption id="Marigold">
[Marigold](https://marigoldmonodepth.github.io/) is a depth estimation diffusion pipeline that uses the rich existing and inherent visual knowledge in diffusion models. It takes an input image and denoises and decodes it into a depth map. Marigold performs well even on images it hasn't seen before.
```py
import torch
from PIL import Image
from diffusers import DiffusionPipeline
from diffusers.utils import load_image
pipeline = DiffusionPipeline.from_pretrained(
"prs-eth/marigold-lcm-v1-0",
custom_pipeline="marigold_depth_estimation",
torch_dtype=torch.float16,
variant="fp16",
)
pipeline.to("cuda")
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/community-marigold.png")
output = pipeline(
image,
denoising_steps=4,
ensemble_size=5,
processing_res=768,
match_input_res=True,
batch_size=0,
seed=33,
color_map="Spectral",
show_progress_bar=True,
)
depth_colored: Image.Image = output.depth_colored
depth_colored.save("./depth_colored.png")
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/community-marigold.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">original image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/marigold-depth.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">colorized depth image</figcaption>
</div>
</div>
</hfoption>
<hfoption id="HD-Painter">
[HD-Painter](https://hf.co/papers/2312.14091) is a high-resolution inpainting pipeline. It introduces a *Prompt-Aware Introverted Attention (PAIntA)* layer to better align a prompt with the area to be inpainted, and *Reweighting Attention Score Guidance (RASG)* to keep the latents more prompt-aligned and within their trained domain to generate realistc images.
```py
import torch
from diffusers import DiffusionPipeline, DDIMScheduler
from diffusers.utils import load_image
pipeline = DiffusionPipeline.from_pretrained(
"Lykon/dreamshaper-8-inpainting",
custom_pipeline="hd_painter"
)
pipeline.scheduler = DDIMScheduler.from_config(pipeline.scheduler.config)
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/hd-painter.jpg")
mask_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/hd-painter-mask.png")
prompt = "football"
image = pipeline(prompt, init_image, mask_image, use_rasg=True, use_painta=True, generator=torch.manual_seed(0)).images[0]
image
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/hd-painter.jpg"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">original image</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/hd-painter-output.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">generated image</figcaption>
</div>
</div>
</hfoption>
</hfoptions>
For more information about community pipelines, take a look at the [Community pipelines](custom_pipeline_examples) guide for how to use them and if you're interested in adding a community pipeline check out the [How to contribute a community pipeline](contribute_pipeline) guide!
## Community components
@@ -266,7 +118,7 @@ Community components allow users to build pipelines that may have customized com
This section shows how users should use community components to build a community pipeline.
You'll use the [showlab/show-1-base](https://huggingface.co/showlab/show-1-base) pipeline checkpoint as an example.
You'll use the [showlab/show-1-base](https://huggingface.co/showlab/show-1-base) pipeline checkpoint as an example. So, let's start loading the components:
1. Import and load the text encoder from Transformers:
@@ -300,17 +152,17 @@ In steps 4 and 5, the custom [UNet](https://github.com/showlab/Show-1/blob/main/
</Tip>
4. Now you'll load a [custom UNet](https://github.com/showlab/Show-1/blob/main/showone/models/unet_3d_condition.py), which in this example, has already been implemented in [showone_unet_3d_condition.py](https://huggingface.co/sayakpaul/show-1-base-with-code/blob/main/unet/showone_unet_3d_condition.py) for your convenience. You'll notice the [`UNet3DConditionModel`] class name is changed to `ShowOneUNet3DConditionModel` because [`UNet3DConditionModel`] already exists in Diffusers. Any components needed for the `ShowOneUNet3DConditionModel` class should be placed in showone_unet_3d_condition.py.
4. Now you'll load a [custom UNet](https://github.com/showlab/Show-1/blob/main/showone/models/unet_3d_condition.py), which in this example, has already been implemented in the `showone_unet_3d_condition.py` [script](https://huggingface.co/sayakpaul/show-1-base-with-code/blob/main/unet/showone_unet_3d_condition.py) for your convenience. You'll notice the `UNet3DConditionModel` class name is changed to `ShowOneUNet3DConditionModel` because [`UNet3DConditionModel`] already exists in Diffusers. Any components needed for the `ShowOneUNet3DConditionModel` class should be placed in the `showone_unet_3d_condition.py` script.
Once this is done, you can initialize the UNet:
Once this is done, you can initialize the UNet:
```python
from showone_unet_3d_condition import ShowOneUNet3DConditionModel
```python
from showone_unet_3d_condition import ShowOneUNet3DConditionModel
unet = ShowOneUNet3DConditionModel.from_pretrained(pipe_id, subfolder="unet")
```
unet = ShowOneUNet3DConditionModel.from_pretrained(pipe_id, subfolder="unet")
```
5. Finally, you'll load the custom pipeline code. For this example, it has already been created for you in [pipeline_t2v_base_pixel.py](https://huggingface.co/sayakpaul/show-1-base-with-code/blob/main/pipeline_t2v_base_pixel.py). This script contains a custom `TextToVideoIFPipeline` class for generating videos from text. Just like the custom UNet, any code needed for the custom pipeline to work should go in pipeline_t2v_base_pixel.py.
5. Finally, you'll load the custom pipeline code. For this example, it has already been created for you in the `pipeline_t2v_base_pixel.py` [script](https://huggingface.co/sayakpaul/show-1-base-with-code/blob/main/pipeline_t2v_base_pixel.py). This script contains a custom `TextToVideoIFPipeline` class for generating videos from text. Just like the custom UNet, any code needed for the custom pipeline to work should go in the `pipeline_t2v_base_pixel.py` script.
Once everything is in place, you can initialize the `TextToVideoIFPipeline` with the `ShowOneUNet3DConditionModel`:
@@ -335,16 +187,13 @@ Push the pipeline to the Hub to share with the community!
pipeline.push_to_hub("custom-t2v-pipeline")
```
After the pipeline is successfully pushed, you need to make a few changes:
After the pipeline is successfully pushed, you need a couple of changes:
1. Change the `_class_name` attribute in [model_index.json](https://huggingface.co/sayakpaul/show-1-base-with-code/blob/main/model_index.json#L2) to `"pipeline_t2v_base_pixel"` and `"TextToVideoIFPipeline"`.
2. Upload `showone_unet_3d_condition.py` to the [unet](https://huggingface.co/sayakpaul/show-1-base-with-code/blob/main/unet/showone_unet_3d_condition.py) subfolder.
3. Upload `pipeline_t2v_base_pixel.py` to the pipeline [repository](https://huggingface.co/sayakpaul/show-1-base-with-code/tree/main).
1. Change the `_class_name` attribute in [`model_index.json`](https://huggingface.co/sayakpaul/show-1-base-with-code/blob/main/model_index.json#L2) to `"pipeline_t2v_base_pixel"` and `"TextToVideoIFPipeline"`.
2. Upload `showone_unet_3d_condition.py` to the `unet` [directory](https://huggingface.co/sayakpaul/show-1-base-with-code/blob/main/unet/showone_unet_3d_condition.py).
3. Upload `pipeline_t2v_base_pixel.py` to the pipeline base [directory](https://huggingface.co/sayakpaul/show-1-base-with-code/blob/main/unet/showone_unet_3d_condition.py).
To run inference, add the `trust_remote_code` argument while initializing the pipeline to handle all the "magic" behind the scenes.
> [!WARNING]
> As an additional precaution with `trust_remote_code=True`, we strongly encourage you to pass a commit hash to the `revision` parameter in [`~DiffusionPipeline.from_pretrained`] to make sure the code hasn't been updated with some malicious new lines of code (unless you fully trust the model owners).
To run inference, simply add the `trust_remote_code` argument while initializing the pipeline to handle all the "magic" behind the scenes.
```python
from diffusers import DiffusionPipeline
@@ -372,9 +221,10 @@ video_frames = pipeline(
).frames
```
As an additional reference, take a look at the repository structure of [stabilityai/japanese-stable-diffusion-xl](https://huggingface.co/stabilityai/japanese-stable-diffusion-xl/) which also uses the `trust_remote_code` feature.
As an additional reference example, you can refer to the repository structure of [stabilityai/japanese-stable-diffusion-xl](https://huggingface.co/stabilityai/japanese-stable-diffusion-xl/), that makes use of the `trust_remote_code` feature:
```python
from diffusers import DiffusionPipeline
import torch
@@ -382,4 +232,12 @@ pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/japanese-stable-diffusion-xl", trust_remote_code=True
)
pipeline.to("cuda")
```
# if using torch < 2.0
# pipeline.enable_xformers_memory_efficient_attention()
prompt = "柴犬、カラフルアート"
image = pipeline(prompt=prompt).images[0]
```

View File

@@ -0,0 +1,135 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Improve generation quality with FreeU
[[open-in-colab]]
The UNet is responsible for denoising during the reverse diffusion process, and there are two distinct features in its architecture:
1. Backbone features primarily contribute to the denoising process
2. Skip features mainly introduce high-frequency features into the decoder module and can make the network overlook the semantics in the backbone features
However, the skip connection can sometimes introduce unnatural image details. [FreeU](https://hf.co/papers/2309.11497) is a technique for improving image quality by rebalancing the contributions from the UNets skip connections and backbone feature maps.
FreeU is applied during inference and it does not require any additional training. The technique works for different tasks such as text-to-image, image-to-image, and text-to-video.
In this guide, you will apply FreeU to the [`StableDiffusionPipeline`], [`StableDiffusionXLPipeline`], and [`TextToVideoSDPipeline`]. You need to install Diffusers from source to run the examples below.
## StableDiffusionPipeline
Load the pipeline:
```py
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, safety_checker=None
).to("cuda")
```
Then enable the FreeU mechanism with the FreeU-specific hyperparameters. These values are scaling factors for the backbone and skip features.
```py
pipeline.enable_freeu(s1=0.9, s2=0.2, b1=1.2, b2=1.4)
```
The values above are from the official FreeU [code repository](https://github.com/ChenyangSi/FreeU) where you can also find [reference hyperparameters](https://github.com/ChenyangSi/FreeU#range-for-more-parameters) for different models.
<Tip>
Disable the FreeU mechanism by calling `disable_freeu()` on a pipeline.
</Tip>
And then run inference:
```py
prompt = "A squirrel eating a burger"
seed = 2023
image = pipeline(prompt, generator=torch.manual_seed(seed)).images[0]
image
```
The figure below compares non-FreeU and FreeU results respectively for the same hyperparameters used above (`prompt` and `seed`):
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/freeu/sdv1_5_freeu.jpg)
Let's see how Stable Diffusion 2 results are impacted:
```py
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-1", torch_dtype=torch.float16, safety_checker=None
).to("cuda")
prompt = "A squirrel eating a burger"
seed = 2023
pipeline.enable_freeu(s1=0.9, s2=0.2, b1=1.1, b2=1.2)
image = pipeline(prompt, generator=torch.manual_seed(seed)).images[0]
image
```
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/freeu/sdv2_1_freeu.jpg)
## Stable Diffusion XL
Finally, let's take a look at how FreeU affects Stable Diffusion XL results:
```py
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16,
).to("cuda")
prompt = "A squirrel eating a burger"
seed = 2023
# Comes from
# https://wandb.ai/nasirk24/UNET-FreeU-SDXL/reports/FreeU-SDXL-Optimal-Parameters--Vmlldzo1NDg4NTUw
pipeline.enable_freeu(s1=0.6, s2=0.4, b1=1.1, b2=1.2)
image = pipeline(prompt, generator=torch.manual_seed(seed)).images[0]
image
```
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/freeu/sdxl_freeu.jpg)
## Text-to-video generation
FreeU can also be used to improve video quality:
```python
from diffusers import DiffusionPipeline
from diffusers.utils import export_to_video
import torch
model_id = "cerspense/zeroscope_v2_576w"
pipe = DiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
prompt = "an astronaut riding a horse on mars"
seed = 2023
# The values come from
# https://github.com/lyn-rgb/FreeU_Diffusers#video-pipelines
pipe.enable_freeu(b1=1.2, b2=1.4, s1=0.9, s2=0.2)
video_frames = pipe(prompt, height=320, width=576, num_frames=30, generator=torch.manual_seed(seed)).frames
export_to_video(video_frames, "astronaut_rides_horse.mp4")
```
Thanks to [kadirnar](https://github.com/kadirnar/) for helping to integrate the feature, and to [justindujardin](https://github.com/justindujardin) for the helpful discussions.

View File

@@ -1,190 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Controlling image quality
The components of a diffusion model, like the UNet and scheduler, can be optimized to improve the quality of generated images leading to better image lighting and details. These techniques are especially useful if you don't have the resources to simply use a larger model for inference. You can enable these techniques during inference without any additional training.
This guide will show you how to turn these techniques on in your pipeline and how to configure them to improve the quality of your generated images.
## Lighting
The Stable Diffusion models aren't very good at generating images that are very bright or dark because the scheduler doesn't start sampling from the last timestep and it doesn't enforce a zero signal-to-noise ratio (SNR). The [Common Diffusion Noise Schedules and Sample Steps are Flawed](https://hf.co/papers/2305.08891) paper fixes these issues which are now available in some Diffusers schedulers.
> [!TIP]
> For inference, you need a model that has been trained with *v_prediction*. To train your own model with *v_prediction*, add the following flag to the [train_text_to_image.py](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image.py) or [train_text_to_image_lora.py](https://github.com/huggingface/diffusers/blob/main/examples/text_to_image/train_text_to_image_lora.py) scripts.
>
> ```bash
> --prediction_type="v_prediction"
> ```
For example, load the [ptx0/pseudo-journey-v2](https://hf.co/ptx0/pseudo-journey-v2) checkpoint which was trained with `v_prediction` and the [`DDIMScheduler`]. Now you should configure the following parameters in the [`DDIMScheduler`].
* `rescale_betas_zero_snr=True` to rescale the noise schedule to zero SNR
* `timestep_spacing="trailing"` to start sampling from the last timestep
Set `guidance_rescale` in the pipeline to prevent over-exposure. A lower value increases brightness but some of the details may appear washed out.
```py
from diffusers import DiffusionPipeline, DDIMScheduler
pipeline = DiffusionPipeline.from_pretrained("ptx0/pseudo-journey-v2", use_safetensors=True)
pipeline.scheduler = DDIMScheduler.from_config(
pipeline.scheduler.config, rescale_betas_zero_snr=True, timestep_spacing="trailing"
)
pipeline.to("cuda")
prompt = "cinematic photo of a snowy mountain at night with the northern lights aurora borealis overhead, 35mm photograph, film, professional, 4k, highly detailed"
generator = torch.Generator(device="cpu").manual_seed(23)
image = pipeline(prompt, guidance_rescale=0.7, generator=generator).images[0]
image
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/no-zero-snr.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">default Stable Diffusion v2-1 image</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/zero-snr.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">image with zero SNR and trailing timestep spacing enabled</figcaption>
</div>
</div>
## Details
[FreeU](https://hf.co/papers/2309.11497) improves image details by rebalancing the UNet's backbone and skip connection weights. The skip connections can cause the model to overlook some of the backbone semantics which may lead to unnatural image details in the generated image. This technique does not require any additional training and can be applied on the fly during inference for tasks like image-to-image and text-to-video.
Use the [`~pipelines.StableDiffusionMixin.enable_freeu`] method on your pipeline and configure the scaling factors for the backbone (`b1` and `b2`) and skip connections (`s1` and `s2`). The number after each scaling factor corresponds to the stage in the UNet where the factor is applied. Take a look at the [FreeU](https://github.com/ChenyangSi/FreeU#parameters) repository for reference hyperparameters for different models.
<hfoptions id="freeu">
<hfoption id="Stable Diffusion v1-5">
```py
import torch
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, safety_checker=None
).to("cuda")
pipeline.enable_freeu(s1=0.9, s2=0.2, b1=1.5, b2=1.6)
generator = torch.Generator(device="cpu").manual_seed(33)
prompt = ""
image = pipeline(prompt, generator=generator).images[0]
image
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/sdv15-no-freeu.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">FreeU disabled</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/sdv15-freeu.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">FreeU enabled</figcaption>
</div>
</div>
</hfoption>
<hfoption id="Stable Diffusion v2-1">
```py
import torch
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-1", torch_dtype=torch.float16, safety_checker=None
).to("cuda")
pipeline.enable_freeu(s1=0.9, s2=0.2, b1=1.4, b2=1.6)
generator = torch.Generator(device="cpu").manual_seed(80)
prompt = "A squirrel eating a burger"
image = pipeline(prompt, generator=generator).images[0]
image
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/sdv21-no-freeu.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">FreeU disabled</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/sdv21-freeu.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">FreeU enabled</figcaption>
</div>
</div>
</hfoption>
<hfoption id="Stable Diffusion XL">
```py
import torch
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16,
).to("cuda")
pipeline.enable_freeu(s1=0.9, s2=0.2, b1=1.3, b2=1.4)
generator = torch.Generator(device="cpu").manual_seed(13)
prompt = "A squirrel eating a burger"
image = pipeline(prompt, generator=generator).images[0]
image
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/sdxl-no-freeu.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">FreeU disabled</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/sdxl-freeu.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">FreeU enabled</figcaption>
</div>
</div>
</hfoption>
<hfoption id="Zeroscope">
```py
import torch
from diffusers import DiffusionPipeline
from diffusers.utils import export_to_video
pipeline = DiffusionPipeline.from_pretrained(
"damo-vilab/text-to-video-ms-1.7b", torch_dtype=torch.float16
).to("cuda")
# values come from https://github.com/lyn-rgb/FreeU_Diffusers#video-pipelines
pipeline.enable_freeu(b1=1.2, b2=1.4, s1=0.9, s2=0.2)
prompt = "Confident teddy bear surfer rides the wave in the tropics"
generator = torch.Generator(device="cpu").manual_seed(47)
video_frames = pipeline(prompt, generator=generator).frames[0]
export_to_video(video_frames, "teddy_bear.mp4", fps=10)
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/video-no-freeu.gif"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">FreeU disabled</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/video-freeu.gif"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">FreeU enabled</figcaption>
</div>
</div>
</hfoption>
</hfoptions>
Call the [`pipelines.StableDiffusionMixin.disable_freeu`] method to disable FreeU.
```py
pipeline.disable_freeu()
```

View File

@@ -1,438 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
[[open-in-colab]]
# Trajectory Consistency Distillation-LoRA
Trajectory Consistency Distillation (TCD) enables a model to generate higher quality and more detailed images with fewer steps. Moreover, owing to the effective error mitigation during the distillation process, TCD demonstrates superior performance even under conditions of large inference steps.
The major advantages of TCD are:
- Better than Teacher: TCD demonstrates superior generative quality at both small and large inference steps and exceeds the performance of [DPM-Solver++(2S)](../../api/schedulers/multistep_dpm_solver) with Stable Diffusion XL (SDXL). There is no additional discriminator or LPIPS supervision included during TCD training.
- Flexible Inference Steps: The inference steps for TCD sampling can be freely adjusted without adversely affecting the image quality.
- Freely change detail level: During inference, the level of detail in the image can be adjusted with a single hyperparameter, *gamma*.
> [!TIP]
> For more technical details of TCD, please refer to the [paper](https://arxiv.org/abs/2402.19159) or official [project page](https://mhh0318.github.io/tcd/)).
For large models like SDXL, TCD is trained with [LoRA](https://huggingface.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora) to reduce memory usage. This is also useful because you can reuse LoRAs between different finetuned models, as long as they share the same base model, without further training.
This guide will show you how to perform inference with TCD-LoRAs for a variety of tasks like text-to-image and inpainting, as well as how you can easily combine TCD-LoRAs with other adapters. Choose one of the supported base model and it's corresponding TCD-LoRA checkpoint from the table below to get started.
| Base model | TCD-LoRA checkpoint |
|-------------------------------------------------------------------------------------------------|----------------------------------------------------------------|
| [stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5) | [TCD-SD15](https://huggingface.co/h1t/TCD-SD15-LoRA) |
| [stable-diffusion-2-1-base](https://huggingface.co/stabilityai/stable-diffusion-2-1-base) | [TCD-SD21-base](https://huggingface.co/h1t/TCD-SD21-base-LoRA) |
| [stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) | [TCD-SDXL](https://huggingface.co/h1t/TCD-SDXL-LoRA) |
Make sure you have [PEFT](https://github.com/huggingface/peft) installed for better LoRA support.
```bash
pip install -U peft
```
## General tasks
In this guide, let's use the [`StableDiffusionXLPipeline`] and the [`TCDScheduler`]. Use the [`~StableDiffusionPipeline.load_lora_weights`] method to load the SDXL-compatible TCD-LoRA weights.
A few tips to keep in mind for TCD-LoRA inference are to:
- Keep the `num_inference_steps` between 4 and 50
- Set `eta` (used to control stochasticity at each step) between 0 and 1. You should use a higher `eta` when increasing the number of inference steps, but the downside is that a larger `eta` in [`TCDScheduler`] leads to blurrier images. A value of 0.3 is recommended to produce good results.
<hfoptions id="tasks">
<hfoption id="text-to-image">
```python
import torch
from diffusers import StableDiffusionXLPipeline, TCDScheduler
device = "cuda"
base_model_id = "stabilityai/stable-diffusion-xl-base-1.0"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
pipe = StableDiffusionXLPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16, variant="fp16").to(device)
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
prompt = "Painting of the orange cat Otto von Garfield, Count of Bismarck-Schönhausen, Duke of Lauenburg, Minister-President of Prussia. Depicted wearing a Prussian Pickelhaube and eating his favorite meal - lasagna."
image = pipe(
prompt=prompt,
num_inference_steps=4,
guidance_scale=0,
eta=0.3,
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/demo_image.png)
</hfoption>
<hfoption id="inpainting">
```python
import torch
from diffusers import AutoPipelineForInpainting, TCDScheduler
from diffusers.utils import load_image, make_image_grid
device = "cuda"
base_model_id = "diffusers/stable-diffusion-xl-1.0-inpainting-0.1"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
pipe = AutoPipelineForInpainting.from_pretrained(base_model_id, torch_dtype=torch.float16, variant="fp16").to(device)
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
init_image = load_image(img_url).resize((1024, 1024))
mask_image = load_image(mask_url).resize((1024, 1024))
prompt = "a tiger sitting on a park bench"
image = pipe(
prompt=prompt,
image=init_image,
mask_image=mask_image,
num_inference_steps=8,
guidance_scale=0,
eta=0.3,
strength=0.99, # make sure to use `strength` below 1.0
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
grid_image = make_image_grid([init_image, mask_image, image], rows=1, cols=3)
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/inpainting_tcd.png)
</hfoption>
</hfoptions>
## Community models
TCD-LoRA also works with many community finetuned models and plugins. For example, load the [animagine-xl-3.0](https://huggingface.co/cagliostrolab/animagine-xl-3.0) checkpoint which is a community finetuned version of SDXL for generating anime images.
```python
import torch
from diffusers import StableDiffusionXLPipeline, TCDScheduler
device = "cuda"
base_model_id = "cagliostrolab/animagine-xl-3.0"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
pipe = StableDiffusionXLPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16, variant="fp16").to(device)
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
prompt = "A man, clad in a meticulously tailored military uniform, stands with unwavering resolve. The uniform boasts intricate details, and his eyes gleam with determination. Strands of vibrant, windswept hair peek out from beneath the brim of his cap."
image = pipe(
prompt=prompt,
num_inference_steps=8,
guidance_scale=0,
eta=0.3,
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/animagine_xl.png)
TCD-LoRA also supports other LoRAs trained on different styles. For example, let's load the [TheLastBen/Papercut_SDXL](https://huggingface.co/TheLastBen/Papercut_SDXL) LoRA and fuse it with the TCD-LoRA with the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] method.
> [!TIP]
> Check out the [Merge LoRAs](merge_loras) guide to learn more about efficient merging methods.
```python
import torch
from diffusers import StableDiffusionXLPipeline
from scheduling_tcd import TCDScheduler
device = "cuda"
base_model_id = "stabilityai/stable-diffusion-xl-base-1.0"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
styled_lora_id = "TheLastBen/Papercut_SDXL"
pipe = StableDiffusionXLPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16, variant="fp16").to(device)
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id, adapter_name="tcd")
pipe.load_lora_weights(styled_lora_id, adapter_name="style")
pipe.set_adapters(["tcd", "style"], adapter_weights=[1.0, 1.0])
prompt = "papercut of a winter mountain, snow"
image = pipe(
prompt=prompt,
num_inference_steps=4,
guidance_scale=0,
eta=0.3,
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/styled_lora.png)
## Adapters
TCD-LoRA is very versatile, and it can be combined with other adapter types like ControlNets, IP-Adapter, and AnimateDiff.
<hfoptions id="adapters">
<hfoption id="ControlNet">
### Depth ControlNet
```python
import torch
import numpy as np
from PIL import Image
from transformers import DPTFeatureExtractor, DPTForDepthEstimation
from diffusers import ControlNetModel, StableDiffusionXLControlNetPipeline
from diffusers.utils import load_image, make_image_grid
from scheduling_tcd import TCDScheduler
device = "cuda"
depth_estimator = DPTForDepthEstimation.from_pretrained("Intel/dpt-hybrid-midas").to(device)
feature_extractor = DPTFeatureExtractor.from_pretrained("Intel/dpt-hybrid-midas")
def get_depth_map(image):
image = feature_extractor(images=image, return_tensors="pt").pixel_values.to(device)
with torch.no_grad(), torch.autocast(device):
depth_map = depth_estimator(image).predicted_depth
depth_map = torch.nn.functional.interpolate(
depth_map.unsqueeze(1),
size=(1024, 1024),
mode="bicubic",
align_corners=False,
)
depth_min = torch.amin(depth_map, dim=[1, 2, 3], keepdim=True)
depth_max = torch.amax(depth_map, dim=[1, 2, 3], keepdim=True)
depth_map = (depth_map - depth_min) / (depth_max - depth_min)
image = torch.cat([depth_map] * 3, dim=1)
image = image.permute(0, 2, 3, 1).cpu().numpy()[0]
image = Image.fromarray((image * 255.0).clip(0, 255).astype(np.uint8))
return image
base_model_id = "stabilityai/stable-diffusion-xl-base-1.0"
controlnet_id = "diffusers/controlnet-depth-sdxl-1.0"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
controlnet = ControlNetModel.from_pretrained(
controlnet_id,
torch_dtype=torch.float16,
variant="fp16",
).to(device)
pipe = StableDiffusionXLControlNetPipeline.from_pretrained(
base_model_id,
controlnet=controlnet,
torch_dtype=torch.float16,
variant="fp16",
).to(device)
pipe.enable_model_cpu_offload()
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
prompt = "stormtrooper lecture, photorealistic"
image = load_image("https://huggingface.co/lllyasviel/sd-controlnet-depth/resolve/main/images/stormtrooper.png")
depth_image = get_depth_map(image)
controlnet_conditioning_scale = 0.5 # recommended for good generalization
image = pipe(
prompt,
image=depth_image,
num_inference_steps=4,
guidance_scale=0,
eta=0.3,
controlnet_conditioning_scale=controlnet_conditioning_scale,
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
grid_image = make_image_grid([depth_image, image], rows=1, cols=2)
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/controlnet_depth_tcd.png)
### Canny ControlNet
```python
import torch
from diffusers import ControlNetModel, StableDiffusionXLControlNetPipeline
from diffusers.utils import load_image, make_image_grid
from scheduling_tcd import TCDScheduler
device = "cuda"
base_model_id = "stabilityai/stable-diffusion-xl-base-1.0"
controlnet_id = "diffusers/controlnet-canny-sdxl-1.0"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
controlnet = ControlNetModel.from_pretrained(
controlnet_id,
torch_dtype=torch.float16,
variant="fp16",
).to(device)
pipe = StableDiffusionXLControlNetPipeline.from_pretrained(
base_model_id,
controlnet=controlnet,
torch_dtype=torch.float16,
variant="fp16",
).to(device)
pipe.enable_model_cpu_offload()
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
prompt = "ultrarealistic shot of a furry blue bird"
canny_image = load_image("https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd_controlnet/bird_canny.png")
controlnet_conditioning_scale = 0.5 # recommended for good generalization
image = pipe(
prompt,
image=canny_image,
num_inference_steps=4,
guidance_scale=0,
eta=0.3,
controlnet_conditioning_scale=controlnet_conditioning_scale,
generator=torch.Generator(device=device).manual_seed(0),
).images[0]
grid_image = make_image_grid([canny_image, image], rows=1, cols=2)
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/controlnet_canny_tcd.png)
<Tip>
The inference parameters in this example might not work for all examples, so we recommend you to try different values for `num_inference_steps`, `guidance_scale`, `controlnet_conditioning_scale` and `cross_attention_kwargs` parameters and choose the best one.
</Tip>
</hfoption>
<hfoption id="IP-Adapter">
This example shows how to use the TCD-LoRA with the [IP-Adapter](https://github.com/tencent-ailab/IP-Adapter/tree/main) and SDXL.
```python
import torch
from diffusers import StableDiffusionXLPipeline
from diffusers.utils import load_image, make_image_grid
from ip_adapter import IPAdapterXL
from scheduling_tcd import TCDScheduler
device = "cuda"
base_model_path = "stabilityai/stable-diffusion-xl-base-1.0"
image_encoder_path = "sdxl_models/image_encoder"
ip_ckpt = "sdxl_models/ip-adapter_sdxl.bin"
tcd_lora_id = "h1t/TCD-SDXL-LoRA"
pipe = StableDiffusionXLPipeline.from_pretrained(
base_model_path,
torch_dtype=torch.float16,
variant="fp16"
)
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
pipe.load_lora_weights(tcd_lora_id)
pipe.fuse_lora()
ip_model = IPAdapterXL(pipe, image_encoder_path, ip_ckpt, device)
ref_image = load_image("https://raw.githubusercontent.com/tencent-ailab/IP-Adapter/main/assets/images/woman.png").resize((512, 512))
prompt = "best quality, high quality, wearing sunglasses"
image = ip_model.generate(
pil_image=ref_image,
prompt=prompt,
scale=0.5,
num_samples=1,
num_inference_steps=4,
guidance_scale=0,
eta=0.3,
seed=0,
)[0]
grid_image = make_image_grid([ref_image, image], rows=1, cols=2)
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/ip_adapter.png)
</hfoption>
<hfoption id="AnimateDiff">
[`AnimateDiff`] allows animating images using Stable Diffusion models. TCD-LoRA can substantially accelerate the process without degrading image quality. The quality of animation with TCD-LoRA and AnimateDiff has a more lucid outcome.
```python
import torch
from diffusers import MotionAdapter, AnimateDiffPipeline, DDIMScheduler
from scheduling_tcd import TCDScheduler
from diffusers.utils import export_to_gif
adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-v1-5")
pipe = AnimateDiffPipeline.from_pretrained(
"frankjoshua/toonyou_beta6",
motion_adapter=adapter,
).to("cuda")
# set TCDScheduler
pipe.scheduler = TCDScheduler.from_config(pipe.scheduler.config)
# load TCD LoRA
pipe.load_lora_weights("h1t/TCD-SD15-LoRA", adapter_name="tcd")
pipe.load_lora_weights("guoyww/animatediff-motion-lora-zoom-in", weight_name="diffusion_pytorch_model.safetensors", adapter_name="motion-lora")
pipe.set_adapters(["tcd", "motion-lora"], adapter_weights=[1.0, 1.2])
prompt = "best quality, masterpiece, 1girl, looking at viewer, blurry background, upper body, contemporary, dress"
generator = torch.manual_seed(0)
frames = pipe(
prompt=prompt,
num_inference_steps=5,
guidance_scale=0,
cross_attention_kwargs={"scale": 1},
num_frames=24,
eta=0.3,
generator=generator
).frames[0]
export_to_gif(frames, "animation.gif")
```
![](https://github.com/jabir-zheng/TCD/raw/main/assets/animation_example.gif)
</hfoption>
</hfoptions>

View File

@@ -25,9 +25,6 @@ Let's take a look at how to use IP-Adapter's image prompting capabilities with t
In all the following examples, you'll see the [`~loaders.IPAdapterMixin.set_ip_adapter_scale`] method. This method controls the amount of text or image conditioning to apply to the model. A value of `1.0` means the model is only conditioned on the image prompt. Lowering this value encourages the model to produce more diverse images, but they may not be as aligned with the image prompt. Typically, a value of `0.5` achieves a good balance between the two prompt types and produces good results.
> [!TIP]
> In the examples below, try adding `low_cpu_mem_usage=True` to the [`~loaders.IPAdapterMixin.load_ip_adapter`] method to speed up the loading time.
<hfoptions id="tasks">
<hfoption id="Text-to-image">
@@ -51,10 +48,10 @@ Create a text prompt and load an image prompt before passing them to the pipelin
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_adapter_diner.png")
generator = torch.Generator(device="cpu").manual_seed(0)
images = pipeline(
prompt="a polar bear sitting in a chair drinking a milkshake",
prompt="a polar bear sitting in a chair drinking a milkshake",
ip_adapter_image=image,
negative_prompt="deformed, ugly, wrong proportion, low res, bad anatomy, worst quality, low quality",
num_inference_steps=100,
num_inference_steps=100,
generator=generator,
).images
images[0]
@@ -234,130 +231,8 @@ export_to_gif(frames, "gummy_bear.gif")
</hfoption>
</hfoptions>
## Configure parameters
There are a couple of IP-Adapter parameters that are useful to know about and can help you with your image generation tasks. These parameters can make your workflow more efficient or give you more control over image generation.
### Image embeddings
IP-Adapter enabled pipelines provide the `ip_adapter_image_embeds` parameter to accept precomputed image embeddings. This is particularly useful in scenarios where you need to run the IP-Adapter pipeline multiple times because you have more than one image. For example, [multi IP-Adapter](#multi-ip-adapter) is a specific use case where you provide multiple styling images to generate a specific image in a specific style. Loading and encoding multiple images each time you use the pipeline would be inefficient. Instead, you can precompute and save the image embeddings to disk (which can save a lot of space if you're using high-quality images) and load them when you need them.
> [!TIP]
> This parameter also gives you the flexibility to load embeddings from other sources. For example, ComfyUI image embeddings for IP-Adapters are compatible with Diffusers and should work ouf-of-the-box!
Call the [`~StableDiffusionPipeline.prepare_ip_adapter_image_embeds`] method to encode and generate the image embeddings. Then you can save them to disk with `torch.save`.
> [!TIP]
> If you're using IP-Adapter with `ip_adapter_image_embedding` instead of `ip_adapter_image`', you can set `load_ip_adapter(image_encoder_folder=None,...)` because you don't need to load an encoder to generate the image embeddings.
```py
image_embeds = pipeline.prepare_ip_adapter_image_embeds(
ip_adapter_image=image,
ip_adapter_image_embeds=None,
device="cuda",
num_images_per_prompt=1,
do_classifier_free_guidance=True,
)
torch.save(image_embeds, "image_embeds.ipadpt")
```
Now load the image embeddings by passing them to the `ip_adapter_image_embeds` parameter.
```py
image_embeds = torch.load("image_embeds.ipadpt")
images = pipeline(
prompt="a polar bear sitting in a chair drinking a milkshake",
ip_adapter_image_embeds=image_embeds,
negative_prompt="deformed, ugly, wrong proportion, low res, bad anatomy, worst quality, low quality",
num_inference_steps=100,
generator=generator,
).images
```
### IP-Adapter masking
Binary masks specify which portion of the output image should be assigned to an IP-Adapter. This is useful for composing more than one IP-Adapter image. For each input IP-Adapter image, you must provide a binary mask.
To start, preprocess the input IP-Adapter images with the [`~image_processor.IPAdapterMaskProcessor.preprocess()`] to generate their masks. For optimal results, provide the output height and width to [`~image_processor.IPAdapterMaskProcessor.preprocess()`]. This ensures masks with different aspect ratios are appropriately stretched. If the input masks already match the aspect ratio of the generated image, you don't have to set the `height` and `width`.
```py
from diffusers.image_processor import IPAdapterMaskProcessor
mask1 = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_mask_mask1.png")
mask2 = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_mask_mask2.png")
output_height = 1024
output_width = 1024
processor = IPAdapterMaskProcessor()
masks = processor.preprocess([mask1, mask2], height=output_height, width=output_width)
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_mask1.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">mask one</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_mask2.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">mask two</figcaption>
</div>
</div>
When there is more than one input IP-Adapter image, load them as a list and provide the IP-Adapter scale list. Each of the input IP-Adapter images here corresponds to one of the masks generated above.
```py
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="sdxl_models", weight_name=["ip-adapter-plus-face_sdxl_vit-h.safetensors"])
pipeline.set_ip_adapter_scale([[0.7, 0.7]]) # one scale for each image-mask pair
face_image1 = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_mask_girl1.png")
face_image2 = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_mask_girl2.png")
ip_images = [[face_image1, face_image2]]
masks = [masks.reshape(1, masks.shape[0], masks.shape[2], masks.shape[3])]
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_girl1.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">IP-Adapter image one</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_girl2.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">IP-Adapter image two</figcaption>
</div>
</div>
Now pass the preprocessed masks to `cross_attention_kwargs` in the pipeline call.
```py
generator = torch.Generator(device="cpu").manual_seed(0)
num_images = 1
image = pipeline(
prompt="2 girls",
ip_adapter_image=ip_images,
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
num_inference_steps=20,
num_images_per_prompt=num_images,
generator=generator,
cross_attention_kwargs={"ip_adapter_masks": masks}
).images[0]
image
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_adapter_attention_mask_result_seed_0.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">IP-Adapter masking applied</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_adapter_no_attention_mask_result_seed_0.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">no IP-Adapter masking applied</figcaption>
</div>
</div>
> While calling `load_ip_adapter()`, pass `low_cpu_mem_usage=True` to speed up the loading time.
## Specific use cases
@@ -365,12 +240,13 @@ IP-Adapter's image prompting and compatibility with other adapters and models ma
### Face model
Generating accurate faces is challenging because they are complex and nuanced. Diffusers supports two IP-Adapter checkpoints specifically trained to generate faces from the [h94/IP-Adapter](https://huggingface.co/h94/IP-Adapter) repository:
Generating accurate faces is challenging because they are complex and nuanced. Diffusers supports two IP-Adapter checkpoints specifically trained to generate faces:
* [ip-adapter-full-face_sd15.safetensors](https://huggingface.co/h94/IP-Adapter/blob/main/models/ip-adapter-full-face_sd15.safetensors) is conditioned with images of cropped faces and removed backgrounds
* [ip-adapter-plus-face_sd15.safetensors](https://huggingface.co/h94/IP-Adapter/blob/main/models/ip-adapter-plus-face_sd15.safetensors) uses patch embeddings and is conditioned with images of cropped faces
Additionally, Diffusers supports all IP-Adapter checkpoints trained with face embeddings extracted by `insightface` face models. Supported models are from the [h94/IP-Adapter-FaceID](https://huggingface.co/h94/IP-Adapter-FaceID) repository.
> [!TIP]
> [IP-Adapter-FaceID](https://huggingface.co/h94/IP-Adapter-FaceID) is a face-specific IP-Adapter trained with face ID embeddings instead of CLIP image embeddings, allowing you to generate more consistent faces in different contexts and styles. Try out this popular [community pipeline](https://github.com/huggingface/diffusers/tree/main/examples/community#ip-adapter-face-id) and see how it compares to the other face IP-Adapters.
For face models, use the [h94/IP-Adapter](https://huggingface.co/h94/IP-Adapter) checkpoint. It is also recommended to use [`DDIMScheduler`] or [`EulerDiscreteScheduler`] for face models.
@@ -394,7 +270,7 @@ generator = torch.Generator(device="cpu").manual_seed(26)
image = pipeline(
prompt="A photo of Einstein as a chef, wearing an apron, cooking in a French restaurant",
ip_adapter_image=image,
negative_prompt="lowres, bad anatomy, worst quality, low quality",
negative_prompt="lowres, bad anatomy, worst quality, low quality",
num_inference_steps=100,
generator=generator,
).images[0]
@@ -412,71 +288,6 @@ image
</div>
</div>
To use IP-Adapter FaceID models, first extract face embeddings with `insightface`. Then pass the list of tensors to the pipeline as `ip_adapter_image_embeds`.
```py
import torch
from diffusers import StableDiffusionPipeline, DDIMScheduler
from diffusers.utils import load_image
from insightface.app import FaceAnalysis
pipeline = StableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
torch_dtype=torch.float16,
).to("cuda")
pipeline.scheduler = DDIMScheduler.from_config(pipeline.scheduler.config)
pipeline.load_ip_adapter("h94/IP-Adapter-FaceID", subfolder=None, weight_name="ip-adapter-faceid_sd15.bin", image_encoder_folder=None)
pipeline.set_ip_adapter_scale(0.6)
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_mask_girl1.png")
ref_images_embeds = []
app = FaceAnalysis(name="buffalo_l", providers=['CUDAExecutionProvider', 'CPUExecutionProvider'])
app.prepare(ctx_id=0, det_size=(640, 640))
image = cv2.cvtColor(np.asarray(image), cv2.COLOR_BGR2RGB)
faces = app.get(image)
image = torch.from_numpy(faces[0].normed_embedding)
ref_images_embeds.append(image.unsqueeze(0))
ref_images_embeds = torch.stack(ref_images_embeds, dim=0).unsqueeze(0)
neg_ref_images_embeds = torch.zeros_like(ref_images_embeds)
id_embeds = torch.cat([neg_ref_images_embeds, ref_images_embeds]).to(dtype=torch.float16, device="cuda")
generator = torch.Generator(device="cpu").manual_seed(42)
images = pipeline(
prompt="A photo of a girl",
ip_adapter_image_embeds=[id_embeds],
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
num_inference_steps=20, num_images_per_prompt=1,
generator=generator
).images
```
Both IP-Adapter FaceID Plus and Plus v2 models require CLIP image embeddings. You can prepare face embeddings as shown previously, then you can extract and pass CLIP embeddings to the hidden image projection layers.
```py
from insightface.utils import face_align
ref_images_embeds = []
ip_adapter_images = []
app = FaceAnalysis(name="buffalo_l", providers=['CUDAExecutionProvider', 'CPUExecutionProvider'])
app.prepare(ctx_id=0, det_size=(640, 640))
image = cv2.cvtColor(np.asarray(image), cv2.COLOR_BGR2RGB)
faces = app.get(image)
ip_adapter_images.append(face_align.norm_crop(image, landmark=faces[0].kps, image_size=224))
image = torch.from_numpy(faces[0].normed_embedding)
ref_images_embeds.append(image.unsqueeze(0))
ref_images_embeds = torch.stack(ref_images_embeds, dim=0).unsqueeze(0)
neg_ref_images_embeds = torch.zeros_like(ref_images_embeds)
id_embeds = torch.cat([neg_ref_images_embeds, ref_images_embeds]).to(dtype=torch.float16, device="cuda")
clip_embeds = pipeline.prepare_ip_adapter_image_embeds(
[ip_adapter_images], None, torch.device("cuda"), num_images, True)[0]
pipeline.unet.encoder_hid_proj.image_projection_layers[0].clip_embeds = clip_embeds.to(dtype=torch.float16)
pipeline.unet.encoder_hid_proj.image_projection_layers[0].shortcut = False # True if Plus v2
```
### Multi IP-Adapter
More than one IP-Adapter can be used at the same time to generate specific images in more diverse styles. For example, you can use IP-Adapter-Face to generate consistent faces and characters, and IP-Adapter Plus to generate those faces in a specific style.
@@ -493,7 +304,7 @@ from transformers import CLIPVisionModelWithProjection
from diffusers.utils import load_image
image_encoder = CLIPVisionModelWithProjection.from_pretrained(
"h94/IP-Adapter",
"h94/IP-Adapter",
subfolder="models/image_encoder",
torch_dtype=torch.float16,
)
@@ -512,8 +323,8 @@ pipeline = AutoPipelineForText2Image.from_pretrained(
)
pipeline.scheduler = DDIMScheduler.from_config(pipeline.scheduler.config)
pipeline.load_ip_adapter(
"h94/IP-Adapter",
subfolder="sdxl_models",
"h94/IP-Adapter",
subfolder="sdxl_models",
weight_name=["ip-adapter-plus_sdxl_vit-h.safetensors", "ip-adapter-plus-face_sdxl_vit-h.safetensors"]
)
pipeline.set_ip_adapter_scale([0.7, 0.3])
@@ -525,7 +336,7 @@ Load an image prompt and a folder containing images of a certain style you want
```py
face_image = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/women_input.png")
style_folder = "https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/style_ziggy"
style_images = [load_image(f"{style_folder}/img{i}.png") for i in range(10)]
style_images = [load_image(f"{style_folder}/img{i}.png") for i in range(10)]
```
<div class="flex flex-row gap-4">
@@ -547,11 +358,10 @@ generator = torch.Generator(device="cpu").manual_seed(0)
image = pipeline(
prompt="wonderwoman",
ip_adapter_image=[style_images, face_image],
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
num_inference_steps=50, num_images_per_prompt=1,
generator=generator,
).images[0]
image
```
<div class="flex justify-center">
@@ -569,14 +379,14 @@ from diffusers import DiffusionPipeline, LCMScheduler
import torch
from diffusers.utils import load_image
model_id = "sd-dreambooth-library/herge-style"
model_id = "sd-dreambooth-library/herge-style"
lcm_lora_id = "latent-consistency/lcm-lora-sdv1-5"
pipeline = DiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter_sd15.bin")
pipeline.load_lora_weights(lcm_lora_id)
pipeline.scheduler = LCMScheduler.from_config(pipeline.scheduler.config)
pipeline.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
pipeline.enable_model_cpu_offload()
```
@@ -645,13 +455,13 @@ Pass the depth map and IP-Adapter image to the pipeline to generate an image.
```py
generator = torch.Generator(device="cpu").manual_seed(33)
image = pipeline(
prompt="best quality, high quality",
prompt="best quality, high quality",
image=depth_map,
ip_adapter_image=ip_adapter_image,
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
num_inference_steps=50,
generator=generator,
).images[0]
).image[0]
image
```
@@ -659,86 +469,82 @@ image
    <img src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ipa-controlnet-out.png" />
</div>
### Style & layout control
### IP-Adapter masking
[InstantStyle](https://arxiv.org/abs/2404.02733) is a plug-and-play method on top of IP-Adapter, which disentangles style and layout from image prompt to control image generation. This way, you can generate images following only the style or layout from image prompt, with significantly improved diversity. This is achieved by only activating IP-Adapters to specific parts of the model.
Binary masks can be used to specify which portion of the output image should be assigned to an IP-Adapter.
For each input IP-Adapter image, a binary mask and an IP-Adapter must be provided.
By default IP-Adapters are inserted to all layers of the model. Use the [`~loaders.IPAdapterMixin.set_ip_adapter_scale`] method with a dictionary to assign scales to IP-Adapter at different layers.
Before passing the masks to the pipeline, it's essential to preprocess them using [`IPAdapterMaskProcessor.preprocess()`].
> [!TIP]
> For optimal results, provide the output height and width to [`IPAdapterMaskProcessor.preprocess()`]. This ensures that masks with differing aspect ratios are appropriately stretched. If the input masks already match the aspect ratio of the generated image, specifying height and width can be omitted.
Here an example with two masks:
```py
from diffusers import AutoPipelineForText2Image
from diffusers.utils import load_image
import torch
from diffusers.image_processor import IPAdapterMaskProcessor
pipeline = AutoPipelineForText2Image.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16).to("cuda")
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="sdxl_models", weight_name="ip-adapter_sdxl.bin")
mask1 = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_mask1.png")
mask2 = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_mask2.png")
scale = {
"down": {"block_2": [0.0, 1.0]},
"up": {"block_0": [0.0, 1.0, 0.0]},
}
pipeline.set_ip_adapter_scale(scale)
```
output_height = 1024
output_width = 1024
This will activate IP-Adapter at the second layer in the model's down-part block 2 and up-part block 0. The former is the layer where IP-Adapter injects layout information and the latter injects style. Inserting IP-Adapter to these two layers you can generate images following both the style and layout from image prompt, but with contents more aligned to text prompt.
```py
style_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/0052a70beed5bf71b92610a43a52df6d286cd5f3/diffusers/rabbit.jpg")
generator = torch.Generator(device="cpu").manual_seed(26)
image = pipeline(
prompt="a cat, masterpiece, best quality, high quality",
ip_adapter_image=style_image,
negative_prompt="text, watermark, lowres, low quality, worst quality, deformed, glitch, low contrast, noisy, saturation, blurry",
guidance_scale=5,
num_inference_steps=30,
generator=generator,
).images[0]
image
processor = IPAdapterMaskProcessor()
masks = processor.preprocess([mask1, mask2], height=output_height, width=output_width)
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/0052a70beed5bf71b92610a43a52df6d286cd5f3/diffusers/rabbit.jpg"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">IP-Adapter image</figcaption>
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_mask1.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">mask one</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/datasets/cat_style_layout.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">generated image</figcaption>
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_mask2.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">mask two</figcaption>
</div>
</div>
In contrast, inserting IP-Adapter to all layers will often generate images that overly focus on image prompt and diminish diversity.
Activate IP-Adapter only in the style layer and then call the pipeline again.
If you have more than one IP-Adapter image, load them into a list, ensuring each image is assigned to a different IP-Adapter.
```py
scale = {
"up": {"block_0": [0.0, 1.0, 0.0]},
}
pipeline.set_ip_adapter_scale(scale)
face_image1 = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_girl1.png")
face_image2 = load_image("https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_girl2.png")
ip_images =[[image1], [image2]]
generator = torch.Generator(device="cpu").manual_seed(26)
image = pipeline(
prompt="a cat, masterpiece, best quality, high quality",
ip_adapter_image=style_image,
negative_prompt="text, watermark, lowres, low quality, worst quality, deformed, glitch, low contrast, noisy, saturation, blurry",
guidance_scale=5,
num_inference_steps=30,
generator=generator,
).images[0]
image
```
<div class="flex flex-row gap-4">
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/datasets/cat_style_only.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">IP-Adapter only in style layer</figcaption>
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_girl1.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">ip adapter image one</figcaption>
</div>
<div class="flex-1">
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/datasets/cat_ip_adapter.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">IP-Adapter in all layers</figcaption>
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/ip_mask_girl2.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">ip adapter image two</figcaption>
</div>
</div>
Note that you don't have to specify all layers in the dictionary. Those not included in the dictionary will be set to scale 0 which means disable IP-Adapter by default.
Pass preprocessed masks to the pipeline using `cross_attention_kwargs` as shown below:
```py
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="sdxl_models", weight_name=["ip-adapter-plus-face_sdxl_vit-h.safetensors"] * 2)
pipeline.set_ip_adapter_scale([0.7] * 2)
generator = torch.Generator(device="cpu").manual_seed(0)
num_images=1
image = pipeline(
prompt="2 girls",
ip_adapter_image=ip_images,
negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
num_inference_steps=20, num_images_per_prompt=num_images,
generator=generator, cross_attention_kwargs={"ip_adapter_masks": masks}
).images[0]
```
<div class="flex justify-center">
    <img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ip_adapter_attention_mask_result_seed_0.png" />
<figcaption class="mt-2 text-center text-sm text-gray-500">output image</figcaption>
</div>

View File

@@ -10,397 +10,316 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Load pipelines
# Load pipelines, models, and schedulers
[[open-in-colab]]
Diffusion systems consist of multiple components like parameterized models and schedulers that interact in complex ways. That is why we designed the [`DiffusionPipeline`] to wrap the complexity of the entire diffusion system into an easy-to-use API. At the same time, the [`DiffusionPipeline`] is entirely customizable so you can modify each component to build a diffusion system for your use case.
Having an easy way to use a diffusion system for inference is essential to 🧨 Diffusers. Diffusion systems often consist of multiple components like parameterized models, tokenizers, and schedulers that interact in complex ways. That is why we designed the [`DiffusionPipeline`] to wrap the complexity of the entire diffusion system into an easy-to-use API, while remaining flexible enough to be adapted for other use cases, such as loading each component individually as building blocks to assemble your own diffusion system.
Everything you need for inference or training is accessible with the `from_pretrained()` method.
This guide will show you how to load:
- pipelines from the Hub and locally
- different components into a pipeline
- multiple pipelines without increasing memory usage
- checkpoint variants such as different floating point types or non-exponential mean averaged (EMA) weights
- models and schedulers
## Load a pipeline
## Diffusion Pipeline
> [!TIP]
> Skip to the [DiffusionPipeline explained](#diffusionpipeline-explained) section if you're interested in an explanation about how the [`DiffusionPipeline`] class works.
<Tip>
There are two ways to load a pipeline for a task:
💡 Skip to the [DiffusionPipeline explained](#diffusionpipeline-explained) section if you are interested in learning in more detail about how the [`DiffusionPipeline`] class works.
1. Load the generic [`DiffusionPipeline`] class and allow it to automatically detect the correct pipeline class from the checkpoint.
2. Load a specific pipeline class for a specific task.
</Tip>
<hfoptions id="pipelines">
<hfoption id="generic pipeline">
The [`DiffusionPipeline`] class is a simple and generic way to load the latest trending diffusion model from the [Hub](https://huggingface.co/models?library=diffusers&sort=trending). It uses the [`~DiffusionPipeline.from_pretrained`] method to automatically detect the correct pipeline class for a task from the checkpoint, downloads and caches all the required configuration and weight files, and returns a pipeline ready for inference.
The [`DiffusionPipeline`] class is the simplest and most generic way to load the latest trending diffusion model from the [Hub](https://huggingface.co/models?library=diffusers&sort=trending). The [`DiffusionPipeline.from_pretrained`] method automatically detects the correct pipeline class from the checkpoint, downloads, and caches all the required configuration and weight files, and returns a pipeline instance ready for inference.
```python
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", use_safetensors=True)
repo_id = "runwayml/stable-diffusion-v1-5"
pipe = DiffusionPipeline.from_pretrained(repo_id, use_safetensors=True)
```
This same checkpoint can also be used for an image-to-image task. The [`DiffusionPipeline`] class can handle any task as long as you provide the appropriate inputs. For example, for an image-to-image task, you need to pass an initial image to the pipeline.
```py
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", use_safetensors=True)
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-init.png")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
image = pipeline("Astronaut in a jungle, cold color palette, muted colors, detailed, 8k", image=init_image).images[0]
```
</hfoption>
<hfoption id="specific pipeline">
Checkpoints can be loaded by their specific pipeline class if you already know it. For example, to load a Stable Diffusion model, use the [`StableDiffusionPipeline`] class.
You can also load a checkpoint with its specific pipeline class. The example above loaded a Stable Diffusion model; to get the same result, use the [`StableDiffusionPipeline`] class:
```python
from diffusers import StableDiffusionPipeline
pipeline = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", use_safetensors=True)
repo_id = "runwayml/stable-diffusion-v1-5"
pipe = StableDiffusionPipeline.from_pretrained(repo_id, use_safetensors=True)
```
This same checkpoint may also be used for another task like image-to-image. To differentiate what task you want to use the checkpoint for, you have to use the corresponding task-specific pipeline class. For example, to use the same checkpoint for image-to-image, use the [`StableDiffusionImg2ImgPipeline`] class.
A checkpoint (such as [`CompVis/stable-diffusion-v1-4`](https://huggingface.co/CompVis/stable-diffusion-v1-4) or [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5)) may also be used for more than one task, like text-to-image or image-to-image. To differentiate what task you want to use the checkpoint for, you have to load it directly with its corresponding task-specific pipeline class:
```py
```python
from diffusers import StableDiffusionImg2ImgPipeline
pipeline = StableDiffusionImg2ImgPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", use_safetensors=True)
repo_id = "runwayml/stable-diffusion-v1-5"
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(repo_id)
```
</hfoption>
</hfoptions>
Use the Space below to gauge a pipeline's memory requirements before you download and load it to see if it runs on your hardware.
<div class="block dark:hidden">
<iframe
src="https://diffusers-compute-pipeline-size.hf.space?__theme=light"
width="850"
height="1600"
></iframe>
</div>
<div class="hidden dark:block">
<iframe
src="https://diffusers-compute-pipeline-size.hf.space?__theme=dark"
width="850"
height="1600"
></iframe>
</div>
### Local pipeline
To load a pipeline locally, use [git-lfs](https://git-lfs.github.com/) to manually download a checkpoint to your local disk.
To load a diffusion pipeline locally, use [`git-lfs`](https://git-lfs.github.com/) to manually download the checkpoint (in this case, [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5)) to your local disk. This creates a local folder, `./stable-diffusion-v1-5`, on your disk:
```bash
git-lfs install
git clone https://huggingface.co/runwayml/stable-diffusion-v1-5
```
This creates a local folder, ./stable-diffusion-v1-5, on your disk and you should pass its path to [`~DiffusionPipeline.from_pretrained`].
Then pass the local path to [`~DiffusionPipeline.from_pretrained`]:
```python
from diffusers import DiffusionPipeline
stable_diffusion = DiffusionPipeline.from_pretrained("./stable-diffusion-v1-5", use_safetensors=True)
repo_id = "./stable-diffusion-v1-5"
stable_diffusion = DiffusionPipeline.from_pretrained(repo_id, use_safetensors=True)
```
The [`~DiffusionPipeline.from_pretrained`] method won't download files from the Hub when it detects a local path, but this also means it won't download and cache the latest changes to a checkpoint.
The [`~DiffusionPipeline.from_pretrained`] method won't download any files from the Hub when it detects a local path, but this also means it won't download and cache the latest changes to a checkpoint.
## Customize a pipeline
### Swap components in a pipeline
You can customize a pipeline by loading different components into it. This is important because you can:
You can customize the default components of any pipeline with another compatible component. Customization is important because:
- change to a scheduler with faster generation speed or higher generation quality depending on your needs (call the `scheduler.compatibles` method on your pipeline to see compatible schedulers)
- change a default pipeline component to a newer and better performing one
- Changing the scheduler is important for exploring the trade-off between generation speed and quality.
- Different components of a model are typically trained independently and you can swap out a component with a better-performing one.
- During finetuning, usually only some components - like the UNet or text encoder - are trained.
For example, let's customize the default [stabilityai/stable-diffusion-xl-base-1.0](https://hf.co/stabilityai/stable-diffusion-xl-base-1.0) checkpoint with:
- The [`HeunDiscreteScheduler`] to generate higher quality images at the expense of slower generation speed. You must pass the `subfolder="scheduler"` parameter in [`~HeunDiscreteScheduler.from_pretrained`] to load the scheduler configuration into the correct [subfolder](https://hf.co/stabilityai/stable-diffusion-xl-base-1.0/tree/main/scheduler) of the pipeline repository.
- A more stable VAE that runs in fp16.
To find out which schedulers are compatible for customization, you can use the `compatibles` method:
```py
from diffusers import StableDiffusionXLPipeline, HeunDiscreteScheduler, AutoencoderKL
import torch
from diffusers import DiffusionPipeline
scheduler = HeunDiscreteScheduler.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", subfolder="scheduler")
vae = AutoencoderKL.from_pretrained("madebyollin/sdxl-vae-fp16-fix", torch_dtype=torch.float16, use_safetensors=True)
repo_id = "runwayml/stable-diffusion-v1-5"
stable_diffusion = DiffusionPipeline.from_pretrained(repo_id, use_safetensors=True)
stable_diffusion.scheduler.compatibles
```
Now pass the new scheduler and VAE to the [`StableDiffusionXLPipeline`].
Let's use the [`SchedulerMixin.from_pretrained`] method to replace the default [`PNDMScheduler`] with a more performant scheduler, [`EulerDiscreteScheduler`]. The `subfolder="scheduler"` argument is required to load the scheduler configuration from the correct [subfolder](https://huggingface.co/runwayml/stable-diffusion-v1-5/tree/main/scheduler) of the pipeline repository.
```py
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
scheduler=scheduler,
vae=vae,
torch_dtype=torch.float16,
variant="fp16",
use_safetensors=True
).to("cuda")
```
## Reuse a pipeline
When you load multiple pipelines that share the same model components, it makes sense to reuse the shared components instead of reloading everything into memory again, especially if your hardware is memory-constrained. For example:
1. You generated an image with the [`StableDiffusionPipeline`] but you want to improve its quality with the [`StableDiffusionSAGPipeline`]. Both of these pipelines share the same pretrained model, so it'd be a waste of memory to load the same model twice.
2. You want to add a model component, like a [`MotionAdapter`](../api/pipelines/animatediff#animatediffpipeline), to [`AnimateDiffPipeline`] which was instantiated from an existing [`StableDiffusionPipeline`]. Again, both pipelines share the same pretrained model, so it'd be a waste of memory to load an entirely new pipeline again.
With the [`DiffusionPipeline.from_pipe`] API, you can switch between multiple pipelines to take advantage of their different features without increasing memory-usage. It is similar to turning on and off a feature in your pipeline.
> [!TIP]
> To switch between tasks (rather than features), use the [`~DiffusionPipeline.from_pipe`] method with the [AutoPipeline](../api/pipelines/auto_pipeline) class, which automatically identifies the pipeline class based on the task (learn more in the [AutoPipeline](../tutorials/autopipeline) tutorial).
Let's start with a [`StableDiffusionPipeline`] and then reuse the loaded model components to create a [`StableDiffusionSAGPipeline`] to increase generation quality. You'll use the [`StableDiffusionPipeline`] with an [IP-Adapter](./ip_adapter) to generate a bear eating pizza.
Then you can pass the new [`EulerDiscreteScheduler`] instance to the `scheduler` argument in [`DiffusionPipeline`]:
```python
from diffusers import DiffusionPipeline, StableDiffusionSAGPipeline
import torch
import gc
from diffusers.utils import load_image
from accelerate.utils import compute_module_sizes
from diffusers import DiffusionPipeline, EulerDiscreteScheduler
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/load_neg_embed.png")
pipe_sd = DiffusionPipeline.from_pretrained("SG161222/Realistic_Vision_V6.0_B1_noVAE", torch_dtype=torch.float16)
pipe_sd.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter_sd15.bin")
pipe_sd.set_ip_adapter_scale(0.6)
pipe_sd.to("cuda")
generator = torch.Generator(device="cpu").manual_seed(33)
out_sd = pipe_sd(
prompt="bear eats pizza",
negative_prompt="wrong white balance, dark, sketches,worst quality,low quality",
ip_adapter_image=image,
num_inference_steps=50,
generator=generator,
).images[0]
out_sd
repo_id = "runwayml/stable-diffusion-v1-5"
scheduler = EulerDiscreteScheduler.from_pretrained(repo_id, subfolder="scheduler")
stable_diffusion = DiffusionPipeline.from_pretrained(repo_id, scheduler=scheduler, use_safetensors=True)
```
<div class="flex justify-center">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/from_pipe_out_sd_0.png"/>
</div>
### Safety checker
For reference, you can check how much memory this process consumed.
```python
def bytes_to_giga_bytes(bytes):
return bytes / 1024 / 1024 / 1024
print(f"Max memory allocated: {bytes_to_giga_bytes(torch.cuda.max_memory_allocated())} GB")
"Max memory allocated: 4.406213283538818 GB"
```
Now, reuse the same pipeline components from [`StableDiffusionPipeline`] in [`StableDiffusionSAGPipeline`] with the [`~DiffusionPipeline.from_pipe`] method.
> [!WARNING]
> Some pipeline methods may not function properly on new pipelines created with [`~DiffusionPipeline.from_pipe`]. For instance, the [`~DiffusionPipeline.enable_model_cpu_offload`] method installs hooks on the model components based on a unique offloading sequence for each pipeline. If the models are executed in a different order in the new pipeline, the CPU offloading may not work correctly.
>
> To ensure everything works as expected, we recommend re-applying a pipeline method on a new pipeline created with [`~DiffusionPipeline.from_pipe`].
```python
pipe_sag = StableDiffusionSAGPipeline.from_pipe(
pipe_sd
)
generator = torch.Generator(device="cpu").manual_seed(33)
out_sag = pipe_sag(
prompt="bear eats pizza",
negative_prompt="wrong white balance, dark, sketches,worst quality,low quality",
ip_adapter_image=image,
num_inference_steps=50,
generator=generator,
guidance_scale=1.0,
sag_scale=0.75
).images[0]
out_sag
```
<div class="flex justify-center">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/from_pipe_out_sag_1.png"/>
</div>
If you check the memory usage, you'll see it remains the same as before because [`StableDiffusionPipeline`] and [`StableDiffusionSAGPipeline`] are sharing the same pipeline components. This allows you to use them interchangeably without any additional memory overhead.
```py
print(f"Max memory allocated: {bytes_to_giga_bytes(torch.cuda.max_memory_allocated())} GB")
"Max memory allocated: 4.406213283538818 GB"
```
Let's animate the image with the [`AnimateDiffPipeline`] and also add a [`MotionAdapter`] module to the pipeline. For the [`AnimateDiffPipeline`], you need to unload the IP-Adapter first and reload it *after* you've created your new pipeline (this only applies to the [`AnimateDiffPipeline`]).
```py
from diffusers import AnimateDiffPipeline, MotionAdapter, DDIMScheduler
from diffusers.utils import export_to_gif
pipe_sag.unload_ip_adapter()
adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-v1-5-2", torch_dtype=torch.float16)
pipe_animate = AnimateDiffPipeline.from_pipe(pipe_sd, motion_adapter=adapter)
pipe_animate.scheduler = DDIMScheduler.from_config(pipe_animate.scheduler.config, beta_schedule="linear")
# load IP-Adapter and LoRA weights again
pipe_animate.load_ip_adapter("h94/IP-Adapter", subfolder="models", weight_name="ip-adapter_sd15.bin")
pipe_animate.load_lora_weights("guoyww/animatediff-motion-lora-zoom-out", adapter_name="zoom-out")
pipe_animate.to("cuda")
generator = torch.Generator(device="cpu").manual_seed(33)
pipe_animate.set_adapters("zoom-out", adapter_weights=0.75)
out = pipe_animate(
prompt="bear eats pizza",
num_frames=16,
num_inference_steps=50,
ip_adapter_image=image,
generator=generator,
).frames[0]
export_to_gif(out, "out_animate.gif")
```
<div class="flex justify-center">
<img class="rounded-xl" src="https://huggingface.co/datasets/YiYiXu/testing-images/resolve/main/from_pipe_out_animate_3.gif"/>
</div>
The [`AnimateDiffPipeline`] is more memory-intensive and consumes 15GB of memory (see the [Memory-usage of from_pipe](#memory-usage-of-from_pipe) section to learn what this means for your memory-usage).
```py
print(f"Max memory allocated: {bytes_to_giga_bytes(torch.cuda.max_memory_allocated())} GB")
"Max memory allocated: 15.178664207458496 GB"
```
### Modify from_pipe components
Pipelines loaded with [`~DiffusionPipeline.from_pipe`] can be customized with different model components or methods. However, whenever you modify the *state* of the model components, it affects all the other pipelines that share the same components. For example, if you call [`~diffusers.loaders.IPAdapterMixin.unload_ip_adapter`] on the [`StableDiffusionSAGPipeline`], you won't be able to use IP-Adapter with the [`StableDiffusionPipeline`] because it's been removed from their shared components.
```py
pipe.sag_unload_ip_adapter()
generator = torch.Generator(device="cpu").manual_seed(33)
out_sd = pipe_sd(
prompt="bear eats pizza",
negative_prompt="wrong white balance, dark, sketches,worst quality,low quality",
ip_adapter_image=image,
num_inference_steps=50,
generator=generator,
).images[0]
"AttributeError: 'NoneType' object has no attribute 'image_projection_layers'"
```
### Memory usage of from_pipe
The memory requirement of loading multiple pipelines with [`~DiffusionPipeline.from_pipe`] is determined by the pipeline with the highest memory-usage regardless of the number of pipelines you create.
| Pipeline | Memory usage (GB) |
|---|---|
| StableDiffusionPipeline | 4.400 |
| StableDiffusionSAGPipeline | 4.400 |
| AnimateDiffPipeline | 15.178 |
The [`AnimateDiffPipeline`] has the highest memory requirement, so the *total memory-usage* is based only on the [`AnimateDiffPipeline`]. Your memory-usage will not increase if you create additional pipelines as long as their memory requirements doesn't exceed that of the [`AnimateDiffPipeline`]. Each pipeline can be used interchangeably without any additional memory overhead.
## Safety checker
Diffusers implements a [safety checker](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/safety_checker.py) for Stable Diffusion models which can generate harmful content. The safety checker screens the generated output against known hardcoded not-safe-for-work (NSFW) content. If for whatever reason you'd like to disable the safety checker, pass `safety_checker=None` to the [`~DiffusionPipeline.from_pretrained`] method.
Diffusion models like Stable Diffusion can generate harmful content, which is why 🧨 Diffusers has a [safety checker](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/stable_diffusion/safety_checker.py) to check generated outputs against known hardcoded NSFW content. If you'd like to disable the safety checker for whatever reason, pass `None` to the `safety_checker` argument:
```python
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", safety_checker=None, use_safetensors=True)
repo_id = "runwayml/stable-diffusion-v1-5"
stable_diffusion = DiffusionPipeline.from_pretrained(repo_id, safety_checker=None, use_safetensors=True)
"""
You have disabled the safety checker for <class 'diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline'> by passing `safety_checker=None`. Ensure that you abide by the conditions of the Stable Diffusion license and do not expose unfiltered results in services or applications open to the public. Both the diffusers team and Hugging Face strongly recommend keeping the safety filter enabled in all public-facing circumstances, disabling it only for use cases that involve analyzing network behavior or auditing its results. For more information, please have a look at https://github.com/huggingface/diffusers/pull/254 .
"""
```
### Reuse components across pipelines
You can also reuse the same components in multiple pipelines to avoid loading the weights into RAM twice. Use the [`~DiffusionPipeline.components`] method to save the components:
```python
from diffusers import StableDiffusionPipeline, StableDiffusionImg2ImgPipeline
model_id = "runwayml/stable-diffusion-v1-5"
stable_diffusion_txt2img = StableDiffusionPipeline.from_pretrained(model_id, use_safetensors=True)
components = stable_diffusion_txt2img.components
```
Then you can pass the `components` to another pipeline without reloading the weights into RAM:
```py
stable_diffusion_img2img = StableDiffusionImg2ImgPipeline(**components)
```
You can also pass the components individually to the pipeline if you want more flexibility over which components to reuse or disable. For example, to reuse the same components in the text-to-image pipeline, except for the safety checker and feature extractor, in the image-to-image pipeline:
```py
from diffusers import StableDiffusionPipeline, StableDiffusionImg2ImgPipeline
model_id = "runwayml/stable-diffusion-v1-5"
stable_diffusion_txt2img = StableDiffusionPipeline.from_pretrained(model_id, use_safetensors=True)
stable_diffusion_img2img = StableDiffusionImg2ImgPipeline(
vae=stable_diffusion_txt2img.vae,
text_encoder=stable_diffusion_txt2img.text_encoder,
tokenizer=stable_diffusion_txt2img.tokenizer,
unet=stable_diffusion_txt2img.unet,
scheduler=stable_diffusion_txt2img.scheduler,
safety_checker=None,
feature_extractor=None,
requires_safety_checker=False,
)
```
## Checkpoint variants
A checkpoint variant is usually a checkpoint whose weights are:
- Stored in a different floating point type, such as [torch.float16](https://pytorch.org/docs/stable/tensors.html#data-types), because it only requires half the bandwidth and storage to download. You can't use this variant if you're continuing training or using a CPU.
- Non-exponential mean averaged (EMA) weights which shouldn't be used for inference. You should use this variant to continue finetuning a model.
- Stored in a different floating point type for lower precision and lower storage, such as [`torch.float16`](https://pytorch.org/docs/stable/tensors.html#data-types), because it only requires half the bandwidth and storage to download. You can't use this variant if you're continuing training or using a CPU.
- Non-exponential mean averaged (EMA) weights, which shouldn't be used for inference. You should use these to continue fine-tuning a model.
> [!TIP]
> When the checkpoints have identical model structures, but they were trained on different datasets and with a different training setup, they should be stored in separate repositories. For example, [stabilityai/stable-diffusion-2](https://hf.co/stabilityai/stable-diffusion-2) and [stabilityai/stable-diffusion-2-1](https://hf.co/stabilityai/stable-diffusion-2-1) are stored in separate repositories.
<Tip>
Otherwise, a variant is **identical** to the original checkpoint. They have exactly the same serialization format (like [safetensors](./using_safetensors)), model structure, and their weights have identical tensor shapes.
💡 When the checkpoints have identical model structures, but they were trained on different datasets and with a different training setup, they should be stored in separate repositories instead of variations (for example, [`stable-diffusion-v1-4`] and [`stable-diffusion-v1-5`]).
| **checkpoint type** | **weight name** | **argument for loading weights** |
|---------------------|---------------------------------------------|----------------------------------|
| original | diffusion_pytorch_model.safetensors | |
| floating point | diffusion_pytorch_model.fp16.safetensors | `variant`, `torch_dtype` |
| non-EMA | diffusion_pytorch_model.non_ema.safetensors | `variant` |
</Tip>
There are two important arguments for loading variants:
Otherwise, a variant is **identical** to the original checkpoint. They have exactly the same serialization format (like [Safetensors](./using_safetensors)), model structure, and weights that have identical tensor shapes.
- `torch_dtype` specifies the floating point precision of the loaded checkpoint. For example, if you want to save bandwidth by loading a fp16 variant, you should set `variant="fp16"` and `torch_dtype=torch.float16` to *convert the weights* to fp16. Otherwise, the fp16 weights are converted to the default fp32 precision.
| **checkpoint type** | **weight name** | **argument for loading weights** |
|---------------------|-------------------------------------|----------------------------------|
| original | diffusion_pytorch_model.bin | |
| floating point | diffusion_pytorch_model.fp16.bin | `variant`, `torch_dtype` |
| non-EMA | diffusion_pytorch_model.non_ema.bin | `variant` |
If you only set `torch_dtype=torch.float16`, the default fp32 weights are downloaded first and then converted to fp16.
There are two important arguments to know for loading variants:
- `variant` specifies which files should be loaded from the repository. For example, if you want to load a non-EMA variant of a UNet from [runwayml/stable-diffusion-v1-5](https://hf.co/runwayml/stable-diffusion-v1-5/tree/main/unet), set `variant="non_ema"` to download the `non_ema` file.
- `torch_dtype` defines the floating point precision of the loaded checkpoints. For example, if you want to save bandwidth by loading a `fp16` variant, you should specify `torch_dtype=torch.float16` to *convert the weights* to `fp16`. Otherwise, the `fp16` weights are converted to the default `fp32` precision. You can also load the original checkpoint without defining the `variant` argument, and convert it to `fp16` with `torch_dtype=torch.float16`. In this case, the default `fp32` weights are downloaded first, and then they're converted to `fp16` after loading.
<hfoptions id="variants">
<hfoption id="fp16">
- `variant` defines which files should be loaded from the repository. For example, if you want to load a `non_ema` variant from the [`diffusers/stable-diffusion-variants`](https://huggingface.co/diffusers/stable-diffusion-variants/tree/main/unet) repository, you should specify `variant="non_ema"` to download the `non_ema` files.
```py
```python
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained(
# load fp16 variant
stable_diffusion = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", variant="fp16", torch_dtype=torch.float16, use_safetensors=True
)
```
</hfoption>
<hfoption id="non-EMA">
```py
pipeline = DiffusionPipeline.from_pretrained(
# load non_ema variant
stable_diffusion = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", variant="non_ema", use_safetensors=True
)
```
</hfoption>
</hfoptions>
Use the `variant` parameter in the [`DiffusionPipeline.save_pretrained`] method to save a checkpoint as a different floating point type or as a non-EMA variant. You should try save a variant to the same folder as the original checkpoint, so you have the option of loading both from the same folder.
<hfoptions id="save">
<hfoption id="fp16">
To save a checkpoint stored in a different floating-point type or as a non-EMA variant, use the [`DiffusionPipeline.save_pretrained`] method and specify the `variant` argument. You should try and save a variant to the same folder as the original checkpoint, so you can load both from the same folder:
```python
from diffusers import DiffusionPipeline
pipeline.save_pretrained("runwayml/stable-diffusion-v1-5", variant="fp16")
# save as fp16 variant
stable_diffusion.save_pretrained("runwayml/stable-diffusion-v1-5", variant="fp16")
# save as non-ema variant
stable_diffusion.save_pretrained("runwayml/stable-diffusion-v1-5", variant="non_ema")
```
</hfoption>
<hfoption id="non_ema">
```py
pipeline.save_pretrained("runwayml/stable-diffusion-v1-5", variant="non_ema")
```
</hfoption>
</hfoptions>
If you don't save the variant to an existing folder, you must specify the `variant` argument otherwise it'll throw an `Exception` because it can't find the original checkpoint.
If you don't save the variant to an existing folder, you must specify the `variant` argument otherwise it'll throw an `Exception` because it can't find the original checkpoint:
```python
# 👎 this won't work
pipeline = DiffusionPipeline.from_pretrained(
stable_diffusion = DiffusionPipeline.from_pretrained(
"./stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
)
# 👍 this works
pipeline = DiffusionPipeline.from_pretrained(
stable_diffusion = DiffusionPipeline.from_pretrained(
"./stable-diffusion-v1-5", variant="fp16", torch_dtype=torch.float16, use_safetensors=True
)
```
<!--
TODO(Patrick) - Make sure to uncomment this part as soon as things are deprecated.
#### Using `revision` to load pipeline variants is deprecated
Previously the `revision` argument of [`DiffusionPipeline.from_pretrained`] was heavily used to
load model variants, e.g.:
```python
from diffusers import DiffusionPipeline
pipe = DiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", revision="fp16", use_safetensors=True)
```
However, this behavior is now deprecated since the "revision" argument should (just as it's done in GitHub) better be used to load model checkpoints from a specific commit or branch in development.
The above example is therefore deprecated and won't be supported anymore for `diffusers >= 1.0.0`.
<Tip warning={true}>
If you load diffusers pipelines or models with `revision="fp16"` or `revision="non_ema"`,
please make sure to update the code and use `variant="fp16"` or `variation="non_ema"` respectively
instead.
</Tip>
-->
## Models
Models are loaded from the [`ModelMixin.from_pretrained`] method, which downloads and caches the latest version of the model weights and configurations. If the latest files are available in the local cache, [`~ModelMixin.from_pretrained`] reuses files in the cache instead of re-downloading them.
Models can be loaded from a subfolder with the `subfolder` argument. For example, the model weights for `runwayml/stable-diffusion-v1-5` are stored in the [`unet`](https://huggingface.co/runwayml/stable-diffusion-v1-5/tree/main/unet) subfolder:
```python
from diffusers import UNet2DConditionModel
repo_id = "runwayml/stable-diffusion-v1-5"
model = UNet2DConditionModel.from_pretrained(repo_id, subfolder="unet", use_safetensors=True)
```
Or directly from a repository's [directory](https://huggingface.co/google/ddpm-cifar10-32/tree/main):
```python
from diffusers import UNet2DModel
repo_id = "google/ddpm-cifar10-32"
model = UNet2DModel.from_pretrained(repo_id, use_safetensors=True)
```
You can also load and save model variants by specifying the `variant` argument in [`ModelMixin.from_pretrained`] and [`ModelMixin.save_pretrained`]:
```python
from diffusers import UNet2DConditionModel
model = UNet2DConditionModel.from_pretrained(
"runwayml/stable-diffusion-v1-5", subfolder="unet", variant="non_ema", use_safetensors=True
)
model.save_pretrained("./local-unet", variant="non_ema")
```
## Schedulers
Schedulers are loaded from the [`SchedulerMixin.from_pretrained`] method, and unlike models, schedulers are **not parameterized** or **trained**; they are defined by a configuration file.
Loading schedulers does not consume any significant amount of memory and the same configuration file can be used for a variety of different schedulers.
For example, the following schedulers are compatible with [`StableDiffusionPipeline`], which means you can load the same scheduler configuration file in any of these classes:
```python
from diffusers import StableDiffusionPipeline
from diffusers import (
DDPMScheduler,
DDIMScheduler,
PNDMScheduler,
LMSDiscreteScheduler,
EulerAncestralDiscreteScheduler,
EulerDiscreteScheduler,
DPMSolverMultistepScheduler,
)
repo_id = "runwayml/stable-diffusion-v1-5"
ddpm = DDPMScheduler.from_pretrained(repo_id, subfolder="scheduler")
ddim = DDIMScheduler.from_pretrained(repo_id, subfolder="scheduler")
pndm = PNDMScheduler.from_pretrained(repo_id, subfolder="scheduler")
lms = LMSDiscreteScheduler.from_pretrained(repo_id, subfolder="scheduler")
euler_anc = EulerAncestralDiscreteScheduler.from_pretrained(repo_id, subfolder="scheduler")
euler = EulerDiscreteScheduler.from_pretrained(repo_id, subfolder="scheduler")
dpm = DPMSolverMultistepScheduler.from_pretrained(repo_id, subfolder="scheduler")
# replace `dpm` with any of `ddpm`, `ddim`, `pndm`, `lms`, `euler_anc`, `euler`
pipeline = StableDiffusionPipeline.from_pretrained(repo_id, scheduler=dpm, use_safetensors=True)
```
## DiffusionPipeline explained
As a class method, [`DiffusionPipeline.from_pretrained`] is responsible for two things:

View File

@@ -103,7 +103,7 @@ image
<Tip>
LoRA is a very general training technique that can be used with other training methods. For example, it is common to train a model with DreamBooth and LoRA. It is also increasingly common to load and merge multiple LoRAs to create new and unique images. You can learn more about it in the in-depth [Merge LoRAs](merge_loras) guide since merging is outside the scope of this loading guide.
LoRA is a very general training technique that can be used with other training methods. For example, it is common to train a model with DreamBooth and LoRA.
</Tip>
@@ -153,51 +153,113 @@ image
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/load_attn_proc.png" />
</div>
<Tip>
For both [`~loaders.LoraLoaderMixin.load_lora_weights`] and [`~loaders.UNet2DConditionLoadersMixin.load_attn_procs`], you can pass the `cross_attention_kwargs={"scale": 0.5}` parameter to adjust how much of the LoRA weights to use. A value of `0` is the same as only using the base model weights, and a value of `1` is equivalent to using the fully finetuned LoRA.
</Tip>
To unload the LoRA weights, use the [`~loaders.LoraLoaderMixin.unload_lora_weights`] method to discard the LoRA weights and restore the model to its original weights:
```py
pipeline.unload_lora_weights()
```
### Adjust LoRA weight scale
### Load multiple LoRAs
For both [`~loaders.LoraLoaderMixin.load_lora_weights`] and [`~loaders.UNet2DConditionLoadersMixin.load_attn_procs`], you can pass the `cross_attention_kwargs={"scale": 0.5}` parameter to adjust how much of the LoRA weights to use. A value of `0` is the same as only using the base model weights, and a value of `1` is equivalent to using the fully finetuned LoRA.
It can be fun to use multiple LoRAs together to create something entirely new and unique. The [`~loaders.LoraLoaderMixin.fuse_lora`] method allows you to fuse the LoRA weights with the original weights of the underlying model.
For more granular control on the amount of LoRA weights used per layer, you can use [`~loaders.LoraLoaderMixin.set_adapters`] and pass a dictionary specifying by how much to scale the weights in each layer by.
```python
pipe = ... # create pipeline
pipe.load_lora_weights(..., adapter_name="my_adapter")
scales = {
"text_encoder": 0.5,
"text_encoder_2": 0.5, # only usable if pipe has a 2nd text encoder
"unet": {
"down": 0.9, # all transformers in the down-part will use scale 0.9
# "mid" # in this example "mid" is not given, therefore all transformers in the mid part will use the default scale 1.0
"up": {
"block_0": 0.6, # all 3 transformers in the 0th block in the up-part will use scale 0.6
"block_1": [0.4, 0.8, 1.0], # the 3 transformers in the 1st block in the up-part will use scales 0.4, 0.8 and 1.0 respectively
}
}
}
pipe.set_adapters("my_adapter", scales)
<Tip>
Fusing the weights can lead to a speedup in inference latency because you don't need to separately load the base model and LoRA! You can save your fused pipeline with [`~DiffusionPipeline.save_pretrained`] to avoid loading and fusing the weights every time you want to use the model.
</Tip>
Load an initial model:
```py
from diffusers import StableDiffusionXLPipeline, AutoencoderKL
import torch
vae = AutoencoderKL.from_pretrained("madebyollin/sdxl-vae-fp16-fix", torch_dtype=torch.float16)
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
vae=vae,
torch_dtype=torch.float16,
).to("cuda")
```
This also works with multiple adapters - see [this guide](https://huggingface.co/docs/diffusers/tutorials/using_peft_for_inference#customize-adapters-strength) for how to do it.
Next, load the LoRA checkpoint and fuse it with the original weights. The `lora_scale` parameter controls how much to scale the output by with the LoRA weights. It is important to make the `lora_scale` adjustments in the [`~loaders.LoraLoaderMixin.fuse_lora`] method because it won't work if you try to pass `scale` to the `cross_attention_kwargs` in the pipeline.
If you need to reset the original model weights for any reason (use a different `lora_scale`), you should use the [`~loaders.LoraLoaderMixin.unfuse_lora`] method.
```py
pipeline.load_lora_weights("ostris/ikea-instructions-lora-sdxl")
pipeline.fuse_lora(lora_scale=0.7)
# to unfuse the LoRA weights
pipeline.unfuse_lora()
```
Then fuse this pipeline with the next set of LoRA weights:
```py
pipeline.load_lora_weights("ostris/super-cereal-sdxl-lora")
pipeline.fuse_lora(lora_scale=0.7)
```
<Tip warning={true}>
Currently, [`~loaders.LoraLoaderMixin.set_adapters`] only supports scaling attention weights. If a LoRA has other parts (e.g., resnets or down-/upsamplers), they will keep a scale of 1.0.
You can't unfuse multiple LoRA checkpoints, so if you need to reset the model to its original weights, you'll need to reload it.
</Tip>
Now you can generate an image that uses the weights from both LoRAs:
```py
prompt = "A cute brown bear eating a slice of pizza, stunning color scheme, masterpiece, illustration"
image = pipeline(prompt).images[0]
image
```
### 🤗 PEFT
<Tip>
Read the [Inference with 🤗 PEFT](../tutorials/using_peft_for_inference) tutorial to learn more about its integration with 🤗 Diffusers and how you can easily work with and juggle multiple adapters. You'll need to install 🤗 Diffusers and PEFT from source to run the example in this section.
</Tip>
Another way you can load and use multiple LoRAs is to specify the `adapter_name` parameter in [`~loaders.LoraLoaderMixin.load_lora_weights`]. This method takes advantage of the 🤗 PEFT integration. For example, load and name both LoRA weights:
```py
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16).to("cuda")
pipeline.load_lora_weights("ostris/ikea-instructions-lora-sdxl", weight_name="ikea_instructions_xl_v1_5.safetensors", adapter_name="ikea")
pipeline.load_lora_weights("ostris/super-cereal-sdxl-lora", weight_name="cereal_box_sdxl_v1.safetensors", adapter_name="cereal")
```
Now use the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] to activate both LoRAs, and you can configure how much weight each LoRA should have on the output:
```py
pipeline.set_adapters(["ikea", "cereal"], adapter_weights=[0.7, 0.5])
```
Then, generate an image:
```py
prompt = "A cute brown bear eating a slice of pizza, stunning color scheme, masterpiece, illustration"
image = pipeline(prompt, num_inference_steps=30, cross_attention_kwargs={"scale": 1.0}).images[0]
image
```
### Kohya and TheLastBen
Other popular LoRA trainers from the community include those by [Kohya](https://github.com/kohya-ss/sd-scripts/) and [TheLastBen](https://github.com/TheLastBen/fast-stable-diffusion). These trainers create different LoRA checkpoints than those trained by 🤗 Diffusers, but they can still be loaded in the same way.
<hfoptions id="other-trainers">
<hfoption id="Kohya">
To load a Kohya LoRA, let's download the [Blueprintify SD XL 1.0](https://civitai.com/models/150986/blueprintify-sd-xl-10) checkpoint from [Civitai](https://civitai.com/) as an example:
Let's download the [Blueprintify SD XL 1.0](https://civitai.com/models/150986/blueprintify-sd-xl-10) checkpoint from [Civitai](https://civitai.com/):
```sh
!wget https://civitai.com/api/download/models/168776 -O blueprintify-sd-xl-10.safetensors
@@ -231,9 +293,6 @@ Some limitations of using Kohya LoRAs with 🤗 Diffusers include:
</Tip>
</hfoption>
<hfoption id="TheLastBen">
Loading a checkpoint from TheLastBen is very similar. For example, to load the [TheLastBen/William_Eggleston_Style_SDXL](https://huggingface.co/TheLastBen/William_Eggleston_Style_SDXL) checkpoint:
```py
@@ -249,9 +308,6 @@ image = pipeline(prompt=prompt).images[0]
image
```
</hfoption>
</hfoptions>
## IP-Adapter
[IP-Adapter](https://ip-adapter.github.io/) is a lightweight adapter that enables image prompting for any diffusion model. This adapter works by decoupling the cross-attention layers of the image and text features. All the other model components are frozen and only the embedded image features in the UNet are trained. As a result, IP-Adapter files are typically only ~100MBs.
@@ -284,9 +340,9 @@ Once loaded, you can use the pipeline with an image and text prompt to guide the
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/load_neg_embed.png")
generator = torch.Generator(device="cpu").manual_seed(33)
images = pipeline(
    prompt='best quality, high quality, wearing sunglasses',
    prompt='best quality, high quality, wearing sunglasses',
    ip_adapter_image=image,
    negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
    negative_prompt="monochrome, lowres, bad anatomy, worst quality, low quality",
    num_inference_steps=50,
    generator=generator,
).images[0]
@@ -299,13 +355,11 @@ images
### IP-Adapter Plus
IP-Adapter relies on an image encoder to generate image features. If the IP-Adapter repository contains an `image_encoder` subfolder, the image encoder is automatically loaded and registered to the pipeline. Otherwise, you'll need to explicitly load the image encoder with a [`~transformers.CLIPVisionModelWithProjection`] model and pass it to the pipeline.
IP-Adapter relies on an image encoder to generate image features. If the IP-Adapter repository contains a `image_encoder` subfolder, the image encoder is automatically loaded and registed to the pipeline. Otherwise, you'll need to explicitly load the image encoder with a [`~transformers.CLIPVisionModelWithProjection`] model and pass it to the pipeline.
This is the case for *IP-Adapter Plus* checkpoints which use the ViT-H image encoder.
```py
from transformers import CLIPVisionModelWithProjection
image_encoder = CLIPVisionModelWithProjection.from_pretrained(
"h94/IP-Adapter",
subfolder="models/image_encoder",
@@ -320,40 +374,3 @@ pipeline = AutoPipelineForText2Image.from_pretrained(
pipeline.load_ip_adapter("h94/IP-Adapter", subfolder="sdxl_models", weight_name="ip-adapter-plus_sdxl_vit-h.safetensors")
```
### IP-Adapter Face ID models
The IP-Adapter FaceID models are experimental IP Adapters that use image embeddings generated by `insightface` instead of CLIP image embeddings. Some of these models also use LoRA to improve ID consistency.
You need to install `insightface` and all its requirements to use these models.
<Tip warning={true}>
As InsightFace pretrained models are available for non-commercial research purposes, IP-Adapter-FaceID models are released exclusively for research purposes and are not intended for commercial use.
</Tip>
```py
pipeline = AutoPipelineForText2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16
).to("cuda")
pipeline.load_ip_adapter("h94/IP-Adapter-FaceID", subfolder=None, weight_name="ip-adapter-faceid_sdxl.bin", image_encoder_folder=None)
```
If you want to use one of the two IP-Adapter FaceID Plus models, you must also load the CLIP image encoder, as this models use both `insightface` and CLIP image embeddings to achieve better photorealism.
```py
from transformers import CLIPVisionModelWithProjection
image_encoder = CLIPVisionModelWithProjection.from_pretrained(
"laion/CLIP-ViT-H-14-laion2B-s32B-b79K",
torch_dtype=torch.float16,
)
pipeline = AutoPipelineForText2Image.from_pretrained(
"runwayml/stable-diffusion-v1-5",
image_encoder=image_encoder,
torch_dtype=torch.float16
).to("cuda")
pipeline.load_ip_adapter("h94/IP-Adapter-FaceID", subfolder=None, weight_name="ip-adapter-faceid-plus_sd15.bin")
```

View File

@@ -0,0 +1,17 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Overview
🧨 Diffusers offers many pipelines, models, and schedulers for generative tasks. To make loading these components as simple as possible, we provide a single and unified method - `from_pretrained()` - that loads any of these components from either the Hugging Face [Hub](https://huggingface.co/models?library=diffusers&sort=downloads) or your local machine. Whenever you load a pipeline or model, the latest files are automatically downloaded and cached so you can quickly reuse them next time without redownloading the files.
This section will show you everything you need to know about loading pipelines, how to load different components in a pipeline, how to load checkpoint variants, and how to load community pipelines. You'll also learn how to load schedulers and compare the speed and quality trade-offs of using different schedulers. Finally, you'll see how to convert and load KerasCV checkpoints so you can use them in PyTorch with 🧨 Diffusers.

View File

@@ -1,266 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Merge LoRAs
It can be fun and creative to use multiple [LoRAs]((https://huggingface.co/docs/peft/conceptual_guides/adapter#low-rank-adaptation-lora)) together to generate something entirely new and unique. This works by merging multiple LoRA weights together to produce images that are a blend of different styles. Diffusers provides a few methods to merge LoRAs depending on *how* you want to merge their weights, which can affect image quality.
This guide will show you how to merge LoRAs using the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] and [`~peft.LoraModel.add_weighted_adapter`] methods. To improve inference speed and reduce memory-usage of merged LoRAs, you'll also see how to use the [`~loaders.LoraLoaderMixin.fuse_lora`] method to fuse the LoRA weights with the original weights of the underlying model.
For this guide, load a Stable Diffusion XL (SDXL) checkpoint and the [KappaNeuro/studio-ghibli-style]() and [Norod78/sdxl-chalkboarddrawing-lora]() LoRAs with the [`~loaders.LoraLoaderMixin.load_lora_weights`] method. You'll need to assign each LoRA an `adapter_name` to combine them later.
```py
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16).to("cuda")
pipeline.load_lora_weights("ostris/ikea-instructions-lora-sdxl", weight_name="ikea_instructions_xl_v1_5.safetensors", adapter_name="ikea")
pipeline.load_lora_weights("lordjia/by-feng-zikai", weight_name="fengzikai_v1.0_XL.safetensors", adapter_name="feng")
```
## set_adapters
The [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] method merges LoRA adapters by concatenating their weighted matrices. Use the adapter name to specify which LoRAs to merge, and the `adapter_weights` parameter to control the scaling for each LoRA. For example, if `adapter_weights=[0.5, 0.5]`, then the merged LoRA output is an average of both LoRAs. Try adjusting the adapter weights to see how it affects the generated image!
```py
pipeline.set_adapters(["ikea", "feng"], adapter_weights=[0.7, 0.8])
generator = torch.manual_seed(0)
prompt = "A bowl of ramen shaped like a cute kawaii bear, by Feng Zikai"
image = pipeline(prompt, generator=generator, cross_attention_kwargs={"scale": 1.0}).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/lora_merge_set_adapters.png"/>
</div>
## add_weighted_adapter
> [!WARNING]
> This is an experimental method that adds PEFTs [`~peft.LoraModel.add_weighted_adapter`] method to Diffusers to enable more efficient merging methods. Check out this [issue](https://github.com/huggingface/diffusers/issues/6892) if you're interested in learning more about the motivation and design behind this integration.
The [`~peft.LoraModel.add_weighted_adapter`] method provides access to more efficient merging method such as [TIES and DARE](https://huggingface.co/docs/peft/developer_guides/model_merging). To use these merging methods, make sure you have the latest stable version of Diffusers and PEFT installed.
```bash
pip install -U diffusers peft
```
There are three steps to merge LoRAs with the [`~peft.LoraModel.add_weighted_adapter`] method:
1. Create a [`~peft.PeftModel`] from the underlying model and LoRA checkpoint.
2. Load a base UNet model and the LoRA adapters.
3. Merge the adapters using the [`~peft.LoraModel.add_weighted_adapter`] method and the merging method of your choice.
Let's dive deeper into what these steps entail.
1. Load a UNet that corresponds to the UNet in the LoRA checkpoint. In this case, both LoRAs use the SDXL UNet as their base model.
```python
from diffusers import UNet2DConditionModel
import torch
unet = UNet2DConditionModel.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
use_safetensors=True,
variant="fp16",
subfolder="unet",
).to("cuda")
```
Load the SDXL pipeline and the LoRA checkpoints, starting with the [ostris/ikea-instructions-lora-sdxl](https://huggingface.co/ostris/ikea-instructions-lora-sdxl) LoRA.
```python
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
variant="fp16",
torch_dtype=torch.float16,
unet=unet
).to("cuda")
pipeline.load_lora_weights("ostris/ikea-instructions-lora-sdxl", weight_name="ikea_instructions_xl_v1_5.safetensors", adapter_name="ikea")
```
Now you'll create a [`~peft.PeftModel`] from the loaded LoRA checkpoint by combining the SDXL UNet and the LoRA UNet from the pipeline.
```python
from peft import get_peft_model, LoraConfig
import copy
sdxl_unet = copy.deepcopy(unet)
ikea_peft_model = get_peft_model(
sdxl_unet,
pipeline.unet.peft_config["ikea"],
adapter_name="ikea"
)
original_state_dict = {f"base_model.model.{k}": v for k, v in pipeline.unet.state_dict().items()}
ikea_peft_model.load_state_dict(original_state_dict, strict=True)
```
> [!TIP]
> You can optionally push the ikea_peft_model to the Hub by calling `ikea_peft_model.push_to_hub("ikea_peft_model", token=TOKEN)`.
Repeat this process to create a [`~peft.PeftModel`] from the [lordjia/by-feng-zikai](https://huggingface.co/lordjia/by-feng-zikai) LoRA.
```python
pipeline.delete_adapters("ikea")
sdxl_unet.delete_adapters("ikea")
pipeline.load_lora_weights("lordjia/by-feng-zikai", weight_name="fengzikai_v1.0_XL.safetensors", adapter_name="feng")
pipeline.set_adapters(adapter_names="feng")
feng_peft_model = get_peft_model(
sdxl_unet,
pipeline.unet.peft_config["feng"],
adapter_name="feng"
)
original_state_dict = {f"base_model.model.{k}": v for k, v in pipe.unet.state_dict().items()}
feng_peft_model.load_state_dict(original_state_dict, strict=True)
```
2. Load a base UNet model and then load the adapters onto it.
```python
from peft import PeftModel
base_unet = UNet2DConditionModel.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
use_safetensors=True,
variant="fp16",
subfolder="unet",
).to("cuda")
model = PeftModel.from_pretrained(base_unet, "stevhliu/ikea_peft_model", use_safetensors=True, subfolder="ikea", adapter_name="ikea")
model.load_adapter("stevhliu/feng_peft_model", use_safetensors=True, subfolder="feng", adapter_name="feng")
```
3. Merge the adapters using the [`~peft.LoraModel.add_weighted_adapter`] method and the merging method of your choice (learn more about other merging methods in this [blog post](https://huggingface.co/blog/peft_merging)). For this example, let's use the `"dare_linear"` method to merge the LoRAs.
> [!WARNING]
> Keep in mind the LoRAs need to have the same rank to be merged!
```python
model.add_weighted_adapter(
adapters=["ikea", "feng"],
weights=[1.0, 1.0],
combination_type="dare_linear",
adapter_name="ikea-feng"
)
model.set_adapters("ikea-feng")
```
Now you can generate an image with the merged LoRA.
```python
model = model.to(dtype=torch.float16, device="cuda")
pipeline = DiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", unet=model, variant="fp16", torch_dtype=torch.float16,
).to("cuda")
image = pipeline("A bowl of ramen shaped like a cute kawaii bear, by Feng Zikai", generator=torch.manual_seed(0)).images[0]
image
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/ikea-feng-dare-linear.png"/>
</div>
## fuse_lora
Both the [`~loaders.UNet2DConditionLoadersMixin.set_adapters`] and [`~peft.LoraModel.add_weighted_adapter`] methods require loading the base model and the LoRA adapters separately which incurs some overhead. The [`~loaders.LoraLoaderMixin.fuse_lora`] method allows you to fuse the LoRA weights directly with the original weights of the underlying model. This way, you're only loading the model once which can increase inference and lower memory-usage.
You can use PEFT to easily fuse/unfuse multiple adapters directly into the model weights (both UNet and text encoder) using the [`~loaders.LoraLoaderMixin.fuse_lora`] method, which can lead to a speed-up in inference and lower VRAM usage.
For example, if you have a base model and adapters loaded and set as active with the following adapter weights:
```py
from diffusers import DiffusionPipeline
import torch
pipeline = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16).to("cuda")
pipeline.load_lora_weights("ostris/ikea-instructions-lora-sdxl", weight_name="ikea_instructions_xl_v1_5.safetensors", adapter_name="ikea")
pipeline.load_lora_weights("lordjia/by-feng-zikai", weight_name="fengzikai_v1.0_XL.safetensors", adapter_name="feng")
pipeline.set_adapters(["ikea", "feng"], adapter_weights=[0.7, 0.8])
```
Fuse these LoRAs into the UNet with the [`~loaders.LoraLoaderMixin.fuse_lora`] method. The `lora_scale` parameter controls how much to scale the output by with the LoRA weights. It is important to make the `lora_scale` adjustments in the [`~loaders.LoraLoaderMixin.fuse_lora`] method because it wont work if you try to pass `scale` to the `cross_attention_kwargs` in the pipeline.
```py
pipeline.fuse_lora(adapter_names=["ikea", "feng"], lora_scale=1.0)
```
Then you should use [`~loaders.LoraLoaderMixin.unload_lora_weights`] to unload the LoRA weights since they've already been fused with the underlying base model. Finally, call [`~DiffusionPipeline.save_pretrained`] to save the fused pipeline locally or you could call [`~DiffusionPipeline.push_to_hub`] to push the fused pipeline to the Hub.
```py
pipeline.unload_lora_weights()
# save locally
pipeline.save_pretrained("path/to/fused-pipeline")
# save to the Hub
pipeline.push_to_hub("fused-ikea-feng")
```
Now you can quickly load the fused pipeline and use it for inference without needing to separately load the LoRA adapters.
```py
pipeline = DiffusionPipeline.from_pretrained(
"username/fused-ikea-feng", torch_dtype=torch.float16,
).to("cuda")
image = pipeline("A bowl of ramen shaped like a cute kawaii bear, by Feng Zikai", generator=torch.manual_seed(0)).images[0]
image
```
You can call [`~loaders.LoraLoaderMixin.unfuse_lora`] to restore the original model's weights (for example, if you want to use a different `lora_scale` value). However, this only works if you've only fused one LoRA adapter to the original model. If you've fused multiple LoRAs, you'll need to reload the model.
```py
pipeline.unfuse_lora()
```
### torch.compile
[torch.compile](../optimization/torch2.0#torchcompile) can speed up your pipeline even more, but the LoRA weights must be fused first and then unloaded. Typically, the UNet is compiled because it is such a computationally intensive component of the pipeline.
```py
from diffusers import DiffusionPipeline
import torch
# load base model and LoRAs
pipeline = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16).to("cuda")
pipeline.load_lora_weights("ostris/ikea-instructions-lora-sdxl", weight_name="ikea_instructions_xl_v1_5.safetensors", adapter_name="ikea")
pipeline.load_lora_weights("lordjia/by-feng-zikai", weight_name="fengzikai_v1.0_XL.safetensors", adapter_name="feng")
# activate both LoRAs and set adapter weights
pipeline.set_adapters(["ikea", "feng"], adapter_weights=[0.7, 0.8])
# fuse LoRAs and unload weights
pipeline.fuse_lora(adapter_names=["ikea", "feng"], lora_scale=1.0)
pipeline.unload_lora_weights()
# torch.compile
pipeline.unet.to(memory_format=torch.channels_last)
pipeline.unet = torch.compile(pipeline.unet, mode="reduce-overhead", fullgraph=True)
image = pipeline("A bowl of ramen shaped like a cute kawaii bear, by Feng Zikai", generator=torch.manual_seed(0)).images[0]
```
Learn more about torch.compile in the [Accelerate inference of text-to-image diffusion models](../tutorials/fast_diffusion#torchcompile) guide.
## Next steps
For more conceptual details about how each merging method works, take a look at the [🤗 PEFT welcomes new merging methods](https://huggingface.co/blog/peft_merging#concatenation-cat) blog post!

View File

@@ -0,0 +1,21 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Using Diffusers with other modalities
Diffusers is in the process of expanding to modalities other than images.
Example type | Colab | Pipeline |
:-------------------------:|:-------------------------:|:-------------------------:|
[Molecule conformation](https://www.nature.com/subjects/molecular-conformation#:~:text=Definition,to%20changes%20in%20their%20environment.) generation | [![Open In Colab](https://colab.research.google.com/assets/colab-badge.svg)](https://colab.research.google.com/github/huggingface/notebooks/blob/main/diffusers/geodiff_molecule_conformation.ipynb) | ❌
More coming soon!

View File

@@ -1,18 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Overview
The inference pipeline supports and enables a wide range of techniques that are divided into two categories:
* Pipeline functionality: these techniques modify the pipeline or extend it for other applications. For example, pipeline callbacks add new features to a pipeline and a pipeline can also be extended for distributed inference.
* Improve inference quality: these techniques increase the visual quality of the generated images. For example, you can enhance your prompts with GPT2 to create better images with lower effort.

View File

@@ -0,0 +1,17 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Overview
A pipeline is an end-to-end class that provides a quick and easy way to use a diffusion system for inference by bundling independently trained models and schedulers together. Certain combinations of models and schedulers define specific pipeline types, like [`StableDiffusionXLPipeline`] or [`StableDiffusionControlNetPipeline`], with specific capabilities. All pipeline types inherit from the base [`DiffusionPipeline`] class; pass it any checkpoint, and it'll automatically detect the pipeline type and load the necessary components.
This section demonstrates how to use specific pipelines such as Stable Diffusion XL, ControlNet, and DiffEdit. You'll also learn how to use a distilled version of the Stable Diffusion model to speed up inference, how to create reproducible pipelines, and how to use and contribute community pipelines.

View File

@@ -0,0 +1,191 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Create reproducible pipelines
[[open-in-colab]]
Reproducibility is important for testing, replicating results, and can even be used to [improve image quality](reusing_seeds). However, the randomness in diffusion models is a desired property because it allows the pipeline to generate different images every time it is run. While you can't expect to get the exact same results across platforms, you can expect results to be reproducible across releases and platforms within a certain tolerance range. Even then, tolerance varies depending on the diffusion pipeline and checkpoint.
This is why it's important to understand how to control sources of randomness in diffusion models or use deterministic algorithms.
<Tip>
💡 We strongly recommend reading PyTorch's [statement about reproducibility](https://pytorch.org/docs/stable/notes/randomness.html):
> Completely reproducible results are not guaranteed across PyTorch releases, individual commits, or different platforms. Furthermore, results may not be reproducible between CPU and GPU executions, even when using identical seeds.
</Tip>
## Control randomness
During inference, pipelines rely heavily on random sampling operations which include creating the
Gaussian noise tensors to denoise and adding noise to the scheduling step.
Take a look at the tensor values in the [`DDIMPipeline`] after two inference steps:
```python
from diffusers import DDIMPipeline
import numpy as np
model_id = "google/ddpm-cifar10-32"
# load model and scheduler
ddim = DDIMPipeline.from_pretrained(model_id, use_safetensors=True)
# run pipeline for just two steps and return numpy tensor
image = ddim(num_inference_steps=2, output_type="np").images
print(np.abs(image).sum())
```
Running the code above prints one value, but if you run it again you get a different value. What is going on here?
Every time the pipeline is run, [`torch.randn`](https://pytorch.org/docs/stable/generated/torch.randn.html) uses a different random seed to create Gaussian noise which is denoised stepwise. This leads to a different result each time it is run, which is great for diffusion pipelines since it generates a different random image each time.
But if you need to reliably generate the same image, that'll depend on whether you're running the pipeline on a CPU or GPU.
### CPU
To generate reproducible results on a CPU, you'll need to use a PyTorch [`Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html) and set a seed:
```python
import torch
from diffusers import DDIMPipeline
import numpy as np
model_id = "google/ddpm-cifar10-32"
# load model and scheduler
ddim = DDIMPipeline.from_pretrained(model_id, use_safetensors=True)
# create a generator for reproducibility
generator = torch.Generator(device="cpu").manual_seed(0)
# run pipeline for just two steps and return numpy tensor
image = ddim(num_inference_steps=2, output_type="np", generator=generator).images
print(np.abs(image).sum())
```
Now when you run the code above, it always prints a value of `1491.1711` no matter what because the `Generator` object with the seed is passed to all the random functions of the pipeline.
If you run this code example on your specific hardware and PyTorch version, you should get a similar, if not the same, result.
<Tip>
💡 It might be a bit unintuitive at first to pass `Generator` objects to the pipeline instead of
just integer values representing the seed, but this is the recommended design when dealing with
probabilistic models in PyTorch, as `Generator`s are *random states* that can be
passed to multiple pipelines in a sequence.
</Tip>
### GPU
Writing a reproducible pipeline on a GPU is a bit trickier, and full reproducibility across different hardware is not guaranteed because matrix multiplication - which diffusion pipelines require a lot of - is less deterministic on a GPU than a CPU. For example, if you run the same code example above on a GPU:
```python
import torch
from diffusers import DDIMPipeline
import numpy as np
model_id = "google/ddpm-cifar10-32"
# load model and scheduler
ddim = DDIMPipeline.from_pretrained(model_id, use_safetensors=True)
ddim.to("cuda")
# create a generator for reproducibility
generator = torch.Generator(device="cuda").manual_seed(0)
# run pipeline for just two steps and return numpy tensor
image = ddim(num_inference_steps=2, output_type="np", generator=generator).images
print(np.abs(image).sum())
```
The result is not the same even though you're using an identical seed because the GPU uses a different random number generator than the CPU.
To circumvent this problem, 🧨 Diffusers has a [`~diffusers.utils.torch_utils.randn_tensor`] function for creating random noise on the CPU, and then moving the tensor to a GPU if necessary. The `randn_tensor` function is used everywhere inside the pipeline, allowing the user to **always** pass a CPU `Generator` even if the pipeline is run on a GPU.
You'll see the results are much closer now!
```python
import torch
from diffusers import DDIMPipeline
import numpy as np
model_id = "google/ddpm-cifar10-32"
# load model and scheduler
ddim = DDIMPipeline.from_pretrained(model_id, use_safetensors=True)
ddim.to("cuda")
# create a generator for reproducibility; notice you don't place it on the GPU!
generator = torch.manual_seed(0)
# run pipeline for just two steps and return numpy tensor
image = ddim(num_inference_steps=2, output_type="np", generator=generator).images
print(np.abs(image).sum())
```
<Tip>
💡 If reproducibility is important, we recommend always passing a CPU generator.
The performance loss is often neglectable, and you'll generate much more similar
values than if the pipeline had been run on a GPU.
</Tip>
Finally, for more complex pipelines such as [`UnCLIPPipeline`], these are often extremely
susceptible to precision error propagation. Don't expect similar results across
different GPU hardware or PyTorch versions. In this case, you'll need to run
exactly the same hardware and PyTorch version for full reproducibility.
## Deterministic algorithms
You can also configure PyTorch to use deterministic algorithms to create a reproducible pipeline. However, you should be aware that deterministic algorithms may be slower than nondeterministic ones and you may observe a decrease in performance. But if reproducibility is important to you, then this is the way to go!
Nondeterministic behavior occurs when operations are launched in more than one CUDA stream. To avoid this, set the environment variable [`CUBLAS_WORKSPACE_CONFIG`](https://docs.nvidia.com/cuda/cublas/index.html#results-reproducibility) to `:16:8` to only use one buffer size during runtime.
PyTorch typically benchmarks multiple algorithms to select the fastest one, but if you want reproducibility, you should disable this feature because the benchmark may select different algorithms each time. Lastly, pass `True` to [`torch.use_deterministic_algorithms`](https://pytorch.org/docs/stable/generated/torch.use_deterministic_algorithms.html) to enable deterministic algorithms.
```py
import os
import torch
os.environ["CUBLAS_WORKSPACE_CONFIG"] = ":16:8"
torch.backends.cudnn.benchmark = False
torch.use_deterministic_algorithms(True)
```
Now when you run the same pipeline twice, you'll get identical results.
```py
import torch
from diffusers import DDIMScheduler, StableDiffusionPipeline
model_id = "runwayml/stable-diffusion-v1-5"
pipe = StableDiffusionPipeline.from_pretrained(model_id, use_safetensors=True).to("cuda")
pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config)
g = torch.Generator(device="cuda")
prompt = "A bear is playing a guitar on Times Square"
g.manual_seed(0)
result1 = pipe(prompt=prompt, num_inference_steps=50, generator=g, output_type="latent").images
g.manual_seed(0)
result2 = pipe(prompt=prompt, num_inference_steps=50, generator=g, output_type="latent").images
print("L_inf dist =", abs(result1 - result2).max())
"L_inf dist = tensor(0., device='cuda:0')"
```

View File

@@ -10,179 +10,72 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Reproducible pipelines
# Improve image quality with deterministic generation
Diffusion models are inherently random which is what allows it to generate different outputs every time it is run. But there are certain times when you want to generate the same output every time, like when you're testing, replicating results, and even [improving image quality](#deterministic-batch-generation). While you can't expect to get identical results across platforms, you can expect reproducible results across releases and platforms within a certain tolerance range (though even this may vary).
[[open-in-colab]]
This guide will show you how to control randomness for deterministic generation on a CPU and GPU.
A common way to improve the quality of generated images is with *deterministic batch generation*, generate a batch of images and select one image to improve with a more detailed prompt in a second round of inference. The key is to pass a list of [`torch.Generator`](https://pytorch.org/docs/stable/generated/torch.Generator.html#generator)'s to the pipeline for batched image generation, and tie each `Generator` to a seed so you can reuse it for an image.
> [!TIP]
> We strongly recommend reading PyTorch's [statement about reproducibility](https://pytorch.org/docs/stable/notes/randomness.html):
>
> "Completely reproducible results are not guaranteed across PyTorch releases, individual commits, or different platforms. Furthermore, results may not be reproducible between CPU and GPU executions, even when using identical seeds."
## Control randomness
During inference, pipelines rely heavily on random sampling operations which include creating the
Gaussian noise tensors to denoise and adding noise to the scheduling step.
Take a look at the tensor values in the [`DDIMPipeline`] after two inference steps.
```python
from diffusers import DDIMPipeline
import numpy as np
ddim = DDIMPipeline.from_pretrained( "google/ddpm-cifar10-32", use_safetensors=True)
image = ddim(num_inference_steps=2, output_type="np").images
print(np.abs(image).sum())
```
Running the code above prints one value, but if you run it again you get a different value.
Each time the pipeline is run, [torch.randn](https://pytorch.org/docs/stable/generated/torch.randn.html) uses a different random seed to create the Gaussian noise tensors. This leads to a different result each time it is run and enables the diffusion pipeline to generate a different random image each time.
But if you need to reliably generate the same image, that depends on whether you're running the pipeline on a CPU or GPU.
> [!TIP]
> It might seem unintuitive to pass `Generator` objects to a pipeline instead of the integer value representing the seed. However, this is the recommended design when working with probabilistic models in PyTorch because a `Generator` is a *random state* that can be passed to multiple pipelines in a sequence. As soon as the `Generator` is consumed, the *state* is changed in place which means even if you passed the same `Generator` to a different pipeline, it won't produce the same result because the state is already changed.
<hfoptions id="hardware">
<hfoption id="CPU">
To generate reproducible results on a CPU, you'll need to use a PyTorch [Generator](https://pytorch.org/docs/stable/generated/torch.Generator.html) and set a seed. Now when you run the code, it always prints a value of `1491.1711` because the `Generator` object with the seed is passed to all the random functions in the pipeline. You should get a similar, if not the same, result on whatever hardware and PyTorch version you're using.
```python
import torch
import numpy as np
from diffusers import DDIMPipeline
ddim = DDIMPipeline.from_pretrained("google/ddpm-cifar10-32", use_safetensors=True)
generator = torch.Generator(device="cpu").manual_seed(0)
image = ddim(num_inference_steps=2, output_type="np", generator=generator).images
print(np.abs(image).sum())
```
</hfoption>
<hfoption id="GPU">
Writing a reproducible pipeline on a GPU is a bit trickier, and full reproducibility across different hardware is not guaranteed because matrix multiplication - which diffusion pipelines require a lot of - is less deterministic on a GPU than a CPU. For example, if you run the same code example from the CPU example, you'll get a different result even though the seed is identical. This is because the GPU uses a different random number generator than the CPU.
```python
import torch
import numpy as np
from diffusers import DDIMPipeline
ddim = DDIMPipeline.from_pretrained("google/ddpm-cifar10-32", use_safetensors=True)
ddim.to("cuda")
generator = torch.Generator(device="cuda").manual_seed(0)
image = ddim(num_inference_steps=2, output_type="np", generator=generator).images
print(np.abs(image).sum())
```
To avoid this issue, Diffusers has a [`~utils.torch_utils.randn_tensor`] function for creating random noise on the CPU, and then moving the tensor to a GPU if necessary. The [`~utils.torch_utils.randn_tensor`] function is used everywhere inside the pipeline. Now you can call [torch.manual_seed](https://pytorch.org/docs/stable/generated/torch.manual_seed.html) which automatically creates a CPU `Generator` that can be passed to the pipeline even if it is being run on a GPU.
```python
import torch
import numpy as np
from diffusers import DDIMPipeline
ddim = DDIMPipeline.from_pretrained("google/ddpm-cifar10-32", use_safetensors=True)
ddim.to("cuda")
generator = torch.manual_seed(0)
image = ddim(num_inference_steps=2, output_type="np", generator=generator).images
print(np.abs(image).sum())
```
> [!TIP]
> If reproducibility is important to your use case, we recommend always passing a CPU `Generator`. The performance loss is often negligible and you'll generate more similar values than if the pipeline had been run on a GPU.
Finally, more complex pipelines such as [`UnCLIPPipeline`], are often extremely
susceptible to precision error propagation. You'll need to use
exactly the same hardware and PyTorch version for full reproducibility.
</hfoption>
</hfoptions>
## Deterministic algorithms
You can also configure PyTorch to use deterministic algorithms to create a reproducible pipeline. The downside is that deterministic algorithms may be slower than non-deterministic ones and you may observe a decrease in performance.
Non-deterministic behavior occurs when operations are launched in more than one CUDA stream. To avoid this, set the environment variable [CUBLAS_WORKSPACE_CONFIG](https://docs.nvidia.com/cuda/cublas/index.html#results-reproducibility) to `:16:8` to only use one buffer size during runtime.
PyTorch typically benchmarks multiple algorithms to select the fastest one, but if you want reproducibility, you should disable this feature because the benchmark may select different algorithms each time. Set Diffusers [enable_full_determinism](https://github.com/huggingface/diffusers/blob/142f353e1c638ff1d20bd798402b68f72c1ebbdd/src/diffusers/utils/testing_utils.py#L861) to enable deterministic algorithms.
Let's use [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5) for example, and generate several versions of the following prompt:
```py
enable_full_determinism()
prompt = "Labrador in the style of Vermeer"
```
Now when you run the same pipeline twice, you'll get identical results.
Instantiate a pipeline with [`DiffusionPipeline.from_pretrained`] and place it on a GPU (if available):
```py
import torch
from diffusers import DDIMScheduler, StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", use_safetensors=True).to("cuda")
pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config)
g = torch.Generator(device="cuda")
prompt = "A bear is playing a guitar on Times Square"
g.manual_seed(0)
result1 = pipe(prompt=prompt, num_inference_steps=50, generator=g, output_type="latent").images
g.manual_seed(0)
result2 = pipe(prompt=prompt, num_inference_steps=50, generator=g, output_type="latent").images
print("L_inf dist =", abs(result1 - result2).max())
"L_inf dist = tensor(0., device='cuda:0')"
```
## Deterministic batch generation
A practical application of creating reproducible pipelines is *deterministic batch generation*. You generate a batch of images and select one image to improve with a more detailed prompt. The main idea is to pass a list of [Generator's](https://pytorch.org/docs/stable/generated/torch.Generator.html) to the pipeline and tie each `Generator` to a seed so you can reuse it.
Let's use the [runwayml/stable-diffusion-v1-5](https://huggingface.co/runwayml/stable-diffusion-v1-5) checkpoint and generate a batch of images.
```py
```python
import torch
from diffusers import DiffusionPipeline
from diffusers.utils import make_image_grid
pipeline = DiffusionPipeline.from_pretrained(
pipe = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
)
pipeline = pipeline.to("cuda")
pipe = pipe.to("cuda")
```
Define four different `Generator`s and assign each `Generator` a seed (`0` to `3`). Then generate a batch of images and pick one to iterate on.
> [!WARNING]
> Use a list comprehension that iterates over the batch size specified in `range()` to create a unique `Generator` object for each image in the batch. If you multiply the `Generator` by the batch size integer, it only creates *one* `Generator` object that is used sequentially for each image in the batch.
>
> ```py
> [torch.Generator().manual_seed(seed)] * 4
> ```
Now, define four different `Generator`s and assign each `Generator` a seed (`0` to `3`) so you can reuse a `Generator` later for a specific image:
```python
generator = [torch.Generator(device="cuda").manual_seed(i) for i in range(4)]
prompt = "Labrador in the style of Vermeer"
images = pipeline(prompt, generator=generator, num_images_per_prompt=4).images[0]
```
<Tip warning={true}>
To create a batched seed, you should use a list comprehension that iterates over the length specified in `range()`. This creates a unique `Generator` object for each image in the batch. If you only multiply the `Generator` by the batch size, this only creates one `Generator` object that is used sequentially for each image in the batch.
For example, if you want to use the same seed to create 4 identical images:
```py
[torch.Generator().manual_seed(seed)] * 4
[torch.Generator().manual_seed(seed) for _ in range(4)]
```
</Tip>
Generate the images and have a look:
```python
images = pipe(prompt, generator=generator, num_images_per_prompt=4).images
make_image_grid(images, rows=2, cols=2)
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/reusabe_seeds.jpg"/>
</div>
![img](https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/reusabe_seeds.jpg)
Let's improve the first image (you can choose any image you want) which corresponds to the `Generator` with seed `0`. Add some additional text to your prompt and then make sure you reuse the same `Generator` with seed `0`. All the generated images should resemble the first image.
In this example, you'll improve upon the first image - but in reality, you can use any image you want (even the image with double sets of eyes!). The first image used the `Generator` with seed `0`, so you'll reuse that `Generator` for the second round of inference. To improve the quality of the image, add some additional text to the prompt:
```python
prompt = [prompt + t for t in [", highly realistic", ", artsy", ", trending", ", colorful"]]
generator = [torch.Generator(device="cuda").manual_seed(0) for i in range(4)]
images = pipeline(prompt, generator=generator).images
```
Create four generators with seed `0`, and generate another batch of images, all of which should look like the first image from the previous round!
```python
images = pipe(prompt, generator=generator).images
make_image_grid(images, rows=2, cols=2)
```
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/reusabe_seeds_2.jpg"/>
</div>
![img](https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/reusabe_seeds_2.jpg)

View File

@@ -10,27 +10,57 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Load schedulers and models
# Schedulers
[[open-in-colab]]
Diffusion pipelines are a collection of interchangeable schedulers and models that can be mixed and matched to tailor a pipeline to a specific use case. The scheduler encapsulates the entire denoising process such as the number of denoising steps and the algorithm for finding the denoised sample. A scheduler is not parameterized or trained so they don't take very much memory. The model is usually only concerned with the forward pass of going from a noisy input to a less noisy sample.
Diffusion pipelines are inherently a collection of diffusion models and schedulers that are partly independent from each other. This means that one is able to switch out parts of the pipeline to better customize
a pipeline to one's use case. The best example of this is the [Schedulers](../api/schedulers/overview).
This guide will show you how to load schedulers and models to customize a pipeline. You'll use the [runwayml/stable-diffusion-v1-5](https://hf.co/runwayml/stable-diffusion-v1-5) checkpoint throughout this guide, so let's load it first.
Whereas diffusion models usually simply define the forward pass from noise to a less noisy sample,
schedulers define the whole denoising process, *i.e.*:
- How many denoising steps?
- Stochastic or deterministic?
- What algorithm to use to find the denoised sample?
```py
import torch
They can be quite complex and often define a trade-off between **denoising speed** and **denoising quality**.
It is extremely difficult to measure quantitatively which scheduler works best for a given diffusion pipeline, so it is often recommended to simply try out which works best.
The following paragraphs show how to do so with the 🧨 Diffusers library.
## Load pipeline
Let's start by loading the [`runwayml/stable-diffusion-v1-5`](https://huggingface.co/runwayml/stable-diffusion-v1-5) model in the [`DiffusionPipeline`]:
```python
from huggingface_hub import login
from diffusers import DiffusionPipeline
import torch
login()
pipeline = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
)
```
You can see what scheduler this pipeline uses with the `pipeline.scheduler` attribute.
Next, we move it to GPU:
```py
```python
pipeline.to("cuda")
```
## Access the scheduler
The scheduler is always one of the components of the pipeline and is usually called `"scheduler"`.
So it can be accessed via the `"scheduler"` property.
```python
pipeline.scheduler
```
**Output**:
```
PNDMScheduler {
"_class_name": "PNDMScheduler",
"_diffusers_version": "0.21.4",
@@ -47,156 +77,235 @@ PNDMScheduler {
}
```
## Load a scheduler
Schedulers are defined by a configuration file that can be used by a variety of schedulers. Load a scheduler with the [`SchedulerMixin.from_pretrained`] method, and specify the `subfolder` parameter to load the configuration file into the correct subfolder of the pipeline repository.
For example, to load the [`DDIMScheduler`]:
```py
from diffusers import DDIMScheduler, DiffusionPipeline
ddim = DDIMScheduler.from_pretrained("runwayml/stable-diffusion-v1-5", subfolder="scheduler")
```
Then you can pass the newly loaded scheduler to the pipeline.
We can see that the scheduler is of type [`PNDMScheduler`].
Cool, now let's compare the scheduler in its performance to other schedulers.
First we define a prompt on which we will test all the different schedulers:
```python
pipeline = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", scheduler=ddim, torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
prompt = "A photograph of an astronaut riding a horse on Mars, high resolution, high definition."
```
Next, we create a generator from a random seed that will ensure that we can generate similar images as well as run the pipeline:
```python
generator = torch.Generator(device="cuda").manual_seed(8)
image = pipeline(prompt, generator=generator).images[0]
image
```
<p align="center">
<br>
<img src="https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/diffusers_docs/astronaut_pndm.png" width="400"/>
<br>
</p>
## Changing the scheduler
Now we show how easy it is to change the scheduler of a pipeline. Every scheduler has a property [`~SchedulerMixin.compatibles`]
which defines all compatible schedulers. You can take a look at all available, compatible schedulers for the Stable Diffusion pipeline as follows.
```python
pipeline.scheduler.compatibles
```
**Output**:
```
[diffusers.utils.dummy_torch_and_torchsde_objects.DPMSolverSDEScheduler,
diffusers.schedulers.scheduling_euler_discrete.EulerDiscreteScheduler,
diffusers.schedulers.scheduling_lms_discrete.LMSDiscreteScheduler,
diffusers.schedulers.scheduling_ddim.DDIMScheduler,
diffusers.schedulers.scheduling_ddpm.DDPMScheduler,
diffusers.schedulers.scheduling_heun_discrete.HeunDiscreteScheduler,
diffusers.schedulers.scheduling_dpmsolver_multistep.DPMSolverMultistepScheduler,
diffusers.schedulers.scheduling_deis_multistep.DEISMultistepScheduler,
diffusers.schedulers.scheduling_pndm.PNDMScheduler,
diffusers.schedulers.scheduling_euler_ancestral_discrete.EulerAncestralDiscreteScheduler,
diffusers.schedulers.scheduling_unipc_multistep.UniPCMultistepScheduler,
diffusers.schedulers.scheduling_k_dpm_2_discrete.KDPM2DiscreteScheduler,
diffusers.schedulers.scheduling_dpmsolver_singlestep.DPMSolverSinglestepScheduler,
diffusers.schedulers.scheduling_k_dpm_2_ancestral_discrete.KDPM2AncestralDiscreteScheduler]
```
Cool, lots of schedulers to look at. Feel free to have a look at their respective class definitions:
- [`EulerDiscreteScheduler`],
- [`LMSDiscreteScheduler`],
- [`DDIMScheduler`],
- [`DDPMScheduler`],
- [`HeunDiscreteScheduler`],
- [`DPMSolverMultistepScheduler`],
- [`DEISMultistepScheduler`],
- [`PNDMScheduler`],
- [`EulerAncestralDiscreteScheduler`],
- [`UniPCMultistepScheduler`],
- [`KDPM2DiscreteScheduler`],
- [`DPMSolverSinglestepScheduler`],
- [`KDPM2AncestralDiscreteScheduler`].
We will now compare the input prompt with all other schedulers. To change the scheduler of the pipeline you can make use of the
convenient [`~ConfigMixin.config`] property in combination with the [`~ConfigMixin.from_config`] function.
```python
pipeline.scheduler.config
```
returns a dictionary of the configuration of the scheduler:
**Output**:
```py
FrozenDict([('num_train_timesteps', 1000),
('beta_start', 0.00085),
('beta_end', 0.012),
('beta_schedule', 'scaled_linear'),
('trained_betas', None),
('skip_prk_steps', True),
('set_alpha_to_one', False),
('prediction_type', 'epsilon'),
('timestep_spacing', 'leading'),
('steps_offset', 1),
('_use_default_values', ['timestep_spacing', 'prediction_type']),
('_class_name', 'PNDMScheduler'),
('_diffusers_version', '0.21.4'),
('clip_sample', False)])
```
This configuration can then be used to instantiate a scheduler
of a different class that is compatible with the pipeline. Here,
we change the scheduler to the [`DDIMScheduler`].
```python
from diffusers import DDIMScheduler
pipeline.scheduler = DDIMScheduler.from_config(pipeline.scheduler.config)
```
Cool, now we can run the pipeline again to compare the generation quality.
```python
generator = torch.Generator(device="cuda").manual_seed(8)
image = pipeline(prompt, generator=generator).images[0]
image
```
<p align="center">
<br>
<img src="https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/diffusers_docs/astronaut_ddim.png" width="400"/>
<br>
</p>
If you are a JAX/Flax user, please check [this section](#changing-the-scheduler-in-flax) instead.
## Compare schedulers
Schedulers have their own unique strengths and weaknesses, making it difficult to quantitatively compare which scheduler works best for a pipeline. You typically have to make a trade-off between denoising speed and denoising quality. We recommend trying out different schedulers to find one that works best for your use case. Call the `pipeline.scheduler.compatibles` attribute to see what schedulers are compatible with a pipeline.
So far we have tried running the stable diffusion pipeline with two schedulers: [`PNDMScheduler`] and [`DDIMScheduler`].
A number of better schedulers have been released that can be run with much fewer steps; let's compare them here:
Let's compare the [`LMSDiscreteScheduler`], [`EulerDiscreteScheduler`], [`EulerAncestralDiscreteScheduler`], and the [`DPMSolverMultistepScheduler`] on the following prompt and seed.
[`LMSDiscreteScheduler`] usually leads to better results:
```py
import torch
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16, use_safetensors=True
).to("cuda")
prompt = "A photograph of an astronaut riding a horse on Mars, high resolution, high definition."
generator = torch.Generator(device="cuda").manual_seed(8)
```
To change the pipelines scheduler, use the [`~ConfigMixin.from_config`] method to load a different scheduler's `pipeline.scheduler.config` into the pipeline.
<hfoptions id="schedulers">
<hfoption id="LMSDiscreteScheduler">
[`LMSDiscreteScheduler`] typically generates higher quality images than the default scheduler.
```py
```python
from diffusers import LMSDiscreteScheduler
pipeline.scheduler = LMSDiscreteScheduler.from_config(pipeline.scheduler.config)
generator = torch.Generator(device="cuda").manual_seed(8)
image = pipeline(prompt, generator=generator).images[0]
image
```
</hfoption>
<hfoption id="EulerDiscreteScheduler">
<p align="center">
<br>
<img src="https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/diffusers_docs/astronaut_lms.png" width="400"/>
<br>
</p>
[`EulerDiscreteScheduler`] can generate higher quality images in just 30 steps.
```py
[`EulerDiscreteScheduler`] and [`EulerAncestralDiscreteScheduler`] can generate high quality results with as little as 30 steps.
```python
from diffusers import EulerDiscreteScheduler
pipeline.scheduler = EulerDiscreteScheduler.from_config(pipeline.scheduler.config)
image = pipeline(prompt, generator=generator).images[0]
generator = torch.Generator(device="cuda").manual_seed(8)
image = pipeline(prompt, generator=generator, num_inference_steps=30).images[0]
image
```
</hfoption>
<hfoption id="EulerAncestralDiscreteScheduler">
<p align="center">
<br>
<img src="https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/diffusers_docs/astronaut_euler_discrete.png" width="400"/>
<br>
</p>
[`EulerAncestralDiscreteScheduler`] can generate higher quality images in just 30 steps.
```py
and:
```python
from diffusers import EulerAncestralDiscreteScheduler
pipeline.scheduler = EulerAncestralDiscreteScheduler.from_config(pipeline.scheduler.config)
image = pipeline(prompt, generator=generator).images[0]
generator = torch.Generator(device="cuda").manual_seed(8)
image = pipeline(prompt, generator=generator, num_inference_steps=30).images[0]
image
```
</hfoption>
<hfoption id="DPMSolverMultistepScheduler">
<p align="center">
<br>
<img src="https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/diffusers_docs/astronaut_euler_ancestral.png" width="400"/>
<br>
</p>
[`DPMSolverMultistepScheduler`] provides a balance between speed and quality and can generate higher quality images in just 20 steps.
```py
[`DPMSolverMultistepScheduler`] gives a reasonable speed/quality trade-off and can be run with as little as 20 steps.
```python
from diffusers import DPMSolverMultistepScheduler
pipeline.scheduler = DPMSolverMultistepScheduler.from_config(pipeline.scheduler.config)
image = pipeline(prompt, generator=generator).images[0]
generator = torch.Generator(device="cuda").manual_seed(8)
image = pipeline(prompt, generator=generator, num_inference_steps=20).images[0]
image
```
</hfoption>
</hfoptions>
<p align="center">
<br>
<img src="https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/diffusers_docs/astronaut_dpm.png" width="400"/>
<br>
</p>
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/diffusers_docs/astronaut_lms.png" />
<figcaption class="mt-2 text-center text-sm text-gray-500">LMSDiscreteScheduler</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/diffusers_docs/astronaut_euler_discrete.png" />
<figcaption class="mt-2 text-center text-sm text-gray-500">EulerDiscreteScheduler</figcaption>
</div>
</div>
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/diffusers_docs/astronaut_euler_ancestral.png" />
<figcaption class="mt-2 text-center text-sm text-gray-500">EulerAncestralDiscreteScheduler</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/diffusers_docs/astronaut_dpm.png" />
<figcaption class="mt-2 text-center text-sm text-gray-500">DPMSolverMultistepScheduler</figcaption>
</div>
</div>
As you can see, most images look very similar and are arguably of very similar quality. It often really depends on the specific use case which scheduler to choose. A good approach is always to run multiple different
schedulers to compare results.
Most images look very similar and are comparable in quality. Again, it often comes down to your specific use case so a good approach is to run multiple different schedulers and compare the results.
## Changing the Scheduler in Flax
### Flax schedulers
If you are a JAX/Flax user, you can also change the default pipeline scheduler. This is a complete example of how to run inference using the Flax Stable Diffusion pipeline and the super-fast [DPM-Solver++ scheduler](../api/schedulers/multistep_dpm_solver):
To compare Flax schedulers, you need to additionally load the scheduler state into the model parameters. For example, let's change the default scheduler in [`FlaxStableDiffusionPipeline`] to use the super fast [`FlaxDPMSolverMultistepScheduler`].
> [!WARNING]
> The [`FlaxLMSDiscreteScheduler`] and [`FlaxDDPMScheduler`] are not compatible with the [`FlaxStableDiffusionPipeline`] yet.
```py
```Python
import jax
import numpy as np
from flax.jax_utils import replicate
from flax.training.common_utils import shard
from diffusers import FlaxStableDiffusionPipeline, FlaxDPMSolverMultistepScheduler
model_id = "runwayml/stable-diffusion-v1-5"
scheduler, scheduler_state = FlaxDPMSolverMultistepScheduler.from_pretrained(
"runwayml/stable-diffusion-v1-5",
model_id,
subfolder="scheduler"
)
pipeline, params = FlaxStableDiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
model_id,
scheduler=scheduler,
revision="bf16",
dtype=jax.numpy.bfloat16,
)
params["scheduler"] = scheduler_state
```
Then you can take advantage of Flax's compatibility with TPUs to generate a number of images in parallel. You'll need to make a copy of the model parameters for each available device and then split the inputs across them to generate your desired number of images.
```py
# Generate 1 image per parallel device (8 on TPUv2-8 or TPUv3-8)
prompt = "A photograph of an astronaut riding a horse on Mars, high resolution, high definition."
prompt = "a photo of an astronaut riding a horse on mars"
num_samples = jax.device_count()
prompt_ids = pipeline.prepare_inputs([prompt] * num_samples)
@@ -212,33 +321,11 @@ images = pipeline(prompt_ids, params, prng_seed, num_inference_steps, jit=True).
images = pipeline.numpy_to_pil(np.asarray(images.reshape((num_samples,) + images.shape[-3:])))
```
## Models
<Tip warning={true}>
Models are loaded from the [`ModelMixin.from_pretrained`] method, which downloads and caches the latest version of the model weights and configurations. If the latest files are available in the local cache, [`~ModelMixin.from_pretrained`] reuses files in the cache instead of re-downloading them.
The following Flax schedulers are _not yet compatible_ with the Flax Stable Diffusion Pipeline:
Models can be loaded from a subfolder with the `subfolder` argument. For example, the model weights for [runwayml/stable-diffusion-v1-5](https://hf.co/runwayml/stable-diffusion-v1-5) are stored in the [unet](https://hf.co/runwayml/stable-diffusion-v1-5/tree/main/unet) subfolder.
- `FlaxLMSDiscreteScheduler`
- `FlaxDDPMScheduler`
```python
from diffusers import UNet2DConditionModel
unet = UNet2DConditionModel.from_pretrained("runwayml/stable-diffusion-v1-5", subfolder="unet", use_safetensors=True)
```
They can also be directly loaded from a [repository](https://huggingface.co/google/ddpm-cifar10-32/tree/main).
```python
from diffusers import UNet2DModel
unet = UNet2DModel.from_pretrained("google/ddpm-cifar10-32", use_safetensors=True)
```
To load and save model variants, specify the `variant` argument in [`ModelMixin.from_pretrained`] and [`ModelMixin.save_pretrained`].
```python
from diffusers import UNet2DConditionModel
unet = UNet2DConditionModel.from_pretrained(
"runwayml/stable-diffusion-v1-5", subfolder="unet", variant="non_ema", use_safetensors=True
)
unet.save_pretrained("./local-unet", variant="non_ema")
```
</Tip>

View File

@@ -63,12 +63,11 @@ from diffusers import StableDiffusionXLPipeline, StableDiffusionXLImg2ImgPipelin
import torch
pipeline = StableDiffusionXLPipeline.from_single_file(
"https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0/blob/main/sd_xl_base_1.0.safetensors",
torch_dtype=torch.float16
"https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0/blob/main/sd_xl_base_1.0.safetensors", torch_dtype=torch.float16, variant="fp16", use_safetensors=True
).to("cuda")
refiner = StableDiffusionXLImg2ImgPipeline.from_single_file(
"https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0/blob/main/sd_xl_refiner_1.0.safetensors", torch_dtype=torch.float16
"https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0/blob/main/sd_xl_refiner_1.0.safetensors", torch_dtype=torch.float16, use_safetensors=True, variant="fp16"
).to("cuda")
```

View File

@@ -31,31 +31,29 @@ Before you begin, make sure you have the following libraries installed:
Model weights may be stored in separate subfolders on the Hub or locally, in which case, you should use the [`~StableDiffusionXLPipeline.from_pretrained`] method:
```py
from diffusers import AutoPipelineForText2Image
from diffusers import AutoPipelineForText2Image, AutoPipelineForImage2Image
import torch
pipeline = AutoPipelineForText2Image.from_pretrained("stabilityai/sdxl-turbo", torch_dtype=torch.float16, variant="fp16")
pipeline = pipeline.to("cuda")
```
You can also use the [`~StableDiffusionXLPipeline.from_single_file`] method to load a model checkpoint stored in a single file format (`.ckpt` or `.safetensors`) from the Hub or locally. For this loading method, you need to set `timestep_spacing="trailing"` (feel free to experiment with the other scheduler config values to get better results):
You can also use the [`~StableDiffusionXLPipeline.from_single_file`] method to load a model checkpoint stored in a single file format (`.ckpt` or `.safetensors`) from the Hub or locally:
```py
from diffusers import StableDiffusionXLPipeline, EulerAncestralDiscreteScheduler
from diffusers import StableDiffusionXLPipeline
import torch
pipeline = StableDiffusionXLPipeline.from_single_file(
"https://huggingface.co/stabilityai/sdxl-turbo/blob/main/sd_xl_turbo_1.0_fp16.safetensors",
torch_dtype=torch.float16, variant="fp16")
"https://huggingface.co/stabilityai/sdxl-turbo/blob/main/sd_xl_turbo_1.0_fp16.safetensors", torch_dtype=torch.float16)
pipeline = pipeline.to("cuda")
pipeline.scheduler = EulerAncestralDiscreteScheduler.from_config(pipeline.scheduler.config, timestep_spacing="trailing")
```
## Text-to-image
For text-to-image, pass a text prompt. By default, SDXL Turbo generates a 512x512 image, and that resolution gives the best results. You can try setting the `height` and `width` parameters to 768x768 or 1024x1024, but you should expect quality degradations when doing so.
Make sure to set `guidance_scale` to 0.0 to disable, as the model was trained without it. A single inference step is enough to generate high quality images.
Make sure to set `guidance_scale` to 0.0 to disable, as the model was trained without it. A single inference step is enough to generate high quality images.
Increasing the number of steps to 2, 3 or 4 should improve image quality.
```py
@@ -77,7 +75,7 @@ image
## Image-to-image
For image-to-image generation, make sure that `num_inference_steps * strength` is larger or equal to 1.
For image-to-image generation, make sure that `num_inference_steps * strength` is larger or equal to 1.
The image-to-image pipeline will run for `int(num_inference_steps * strength)` steps, e.g. `0.5 * 2.0 = 1` step in
our example below.
@@ -86,14 +84,14 @@ from diffusers import AutoPipelineForImage2Image
from diffusers.utils import load_image, make_image_grid
# use from_pipe to avoid consuming additional memory when loading a checkpoint
pipeline_image2image = AutoPipelineForImage2Image.from_pipe(pipeline_text2image).to("cuda")
pipeline = AutoPipelineForImage2Image.from_pipe(pipeline_text2image).to("cuda")
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/cat.png")
init_image = init_image.resize((512, 512))
prompt = "cat wizard, gandalf, lord of the rings, detailed, fantasy, cute, adorable, Pixar, Disney, 8k"
image = pipeline_image2image(prompt, image=init_image, strength=0.5, guidance_scale=0.0, num_inference_steps=2).images[0]
image = pipeline(prompt, image=init_image, strength=0.5, guidance_scale=0.0, num_inference_steps=2).images[0]
make_image_grid([init_image, image], rows=1, cols=2)
```
@@ -103,7 +101,7 @@ make_image_grid([init_image, image], rows=1, cols=2)
## Speed-up SDXL Turbo even more
- Compile the UNet if you are using PyTorch version 2.0 or higher. The first inference run will be very slow, but subsequent ones will be much faster.
- Compile the UNet if you are using PyTorch version 2 or better. The first inference run will be very slow, but subsequent ones will be much faster.
```py
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)

View File

@@ -217,9 +217,3 @@ Check your image dimensions to see if they're correct:
images.shape
# (8, 1, 512, 512, 3)
```
## Resources
To learn more about how JAX works with Stable Diffusion, you may be interested in reading:
* [Accelerating Stable Diffusion XL Inference with JAX on Cloud TPU v5e](https://hf.co/blog/sdxl_jax)

View File

@@ -21,7 +21,7 @@ This guide will show you how to use SVD to generate short videos from images.
Before you begin, make sure you have the following libraries installed:
```py
!pip install -q -U diffusers transformers accelerate
!pip install -q -U diffusers transformers accelerate
```
The are two variants of this model, [SVD](https://huggingface.co/stabilityai/stable-video-diffusion-img2vid) and [SVD-XT](https://huggingface.co/stabilityai/stable-video-diffusion-img2vid-xt). The SVD checkpoint is trained to generate 14 frames and the SVD-XT checkpoint is further finetuned to generate 25 frames.
@@ -86,7 +86,7 @@ Video generation is very memory intensive because you're essentially generating
+ frames = pipe(image, decode_chunk_size=2, generator=generator, num_frames=25).frames[0]
```
Using all these tricks together should lower the memory requirement to less than 8GB VRAM.
Using all these tricks togethere should lower the memory requirement to less than 8GB VRAM.
## Micro-conditioning

Some files were not shown because too many files have changed in this diff Show More