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4 Commits

Author SHA1 Message Date
DN6
4838b5b884 Merge branch 'main' into torch-regression-tests 2025-01-07 11:07:42 +05:30
DN6
240ba165fd Merge branch 'main' into torch-regression-tests 2025-01-07 11:04:10 +05:30
DN6
64ae481eb8 update 2025-01-07 08:59:19 +05:30
DN6
be75474ad9 update 2025-01-07 08:43:50 +05:30
1120 changed files with 18001 additions and 120392 deletions

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@@ -1,38 +0,0 @@
name: "\U0001F31F Remote VAE"
description: Feedback for remote VAE pilot
labels: [ "Remote VAE" ]
body:
- type: textarea
id: positive
validations:
required: true
attributes:
label: Did you like the remote VAE solution?
description: |
If you liked it, we would appreciate it if you could elaborate what you liked.
- type: textarea
id: feedback
validations:
required: true
attributes:
label: What can be improved about the current solution?
description: |
Let us know the things you would like to see improved. Note that we will work optimizing the solution once the pilot is over and we have usage.
- type: textarea
id: others
validations:
required: true
attributes:
label: What other VAEs you would like to see if the pilot goes well?
description: |
Provide a list of the VAEs you would like to see in the future if the pilot goes well.
- type: textarea
id: additional-info
attributes:
label: Notify the members of the team
description: |
Tag the following folks when submitting this feedback: @hlky @sayakpaul

View File

@@ -38,7 +38,6 @@ jobs:
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install pandas peft
python -m uv pip uninstall transformers && python -m uv pip install transformers==4.48.0
- name: Environment
run: |
python utils/print_env.py

View File

@@ -142,7 +142,6 @@ jobs:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
RUN_COMPILE: yes
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
@@ -181,55 +180,6 @@ jobs:
pip install slack_sdk tabulate
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
run_torch_compile_tests:
name: PyTorch Compile CUDA tests
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-compile-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test,training]
- name: Environment
run: |
python utils/print_env.py
- name: Run torch compile tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
RUN_COMPILE: yes
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "compile" --make-reports=tests_torch_compile_cuda tests/
- name: Failure short reports
if: ${{ failure() }}
run: cat reports/tests_torch_compile_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: torch_compile_test_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
run_big_gpu_torch_tests:
name: Torch tests on big GPU
strategy:
@@ -315,14 +265,14 @@ jobs:
- name: Run PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_minimum_version_cuda \
tests/models/test_modeling_common.py \
tests/models/test_modelling_common.py \
tests/pipelines/test_pipelines_common.py \
tests/pipelines/test_pipeline_utils.py \
tests/pipelines/test_pipelines.py \
@@ -464,16 +414,10 @@ jobs:
config:
- backend: "bitsandbytes"
test_location: "bnb"
additional_deps: ["peft"]
- backend: "gguf"
test_location: "gguf"
additional_deps: ["peft"]
- backend: "torchao"
test_location: "torchao"
additional_deps: []
- backend: "optimum_quanto"
test_location: "quanto"
additional_deps: []
runs-on:
group: aws-g6e-xlarge-plus
container:
@@ -491,9 +435,6 @@ jobs:
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install -U ${{ matrix.config.backend }}
if [ "${{ join(matrix.config.additional_deps, ' ') }}" != "" ]; then
python -m uv pip install ${{ join(matrix.config.additional_deps, ' ') }}
fi
python -m uv pip install pytest-reportlog
- name: Environment
run: |
@@ -526,60 +467,6 @@ jobs:
pip install slack_sdk tabulate
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
run_nightly_pipeline_level_quantization_tests:
name: Torch quantization nightly tests
strategy:
fail-fast: false
max-parallel: 2
runs-on:
group: aws-g6e-xlarge-plus
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "20gb" --ipc host --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install -U bitsandbytes optimum_quanto
python -m uv pip install pytest-reportlog
- name: Environment
run: |
python utils/print_env.py
- name: Pipeline-level quantization tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
BIG_GPU_MEMORY: 40
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
--make-reports=tests_pipeline_level_quant_torch_cuda \
--report-log=tests_pipeline_level_quant_torch_cuda.log \
tests/quantization/test_pipeline_level_quantization.py
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_level_quant_torch_cuda_stats.txt
cat reports/tests_pipeline_level_quant_torch_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: torch_cuda_pipeline_level_quant_reports
path: reports
- name: Generate Report and Notify Channel
if: always()
run: |
pip install slack_sdk tabulate
python utils/log_reports.py >> $GITHUB_STEP_SUMMARY
# M1 runner currently not well supported
# TODO: (Dhruv) add these back when we setup better testing for Apple Silicon
# run_nightly_tests_apple_m1:
@@ -618,7 +505,7 @@ jobs:
# shell: arch -arch arm64 bash {0}
# env:
# HF_HOME: /System/Volumes/Data/mnt/cache
# HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# HF_TOKEN: ${{ secrets.HF_TOKEN }}
# run: |
# ${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
# --report-log=tests_torch_mps.log \
@@ -674,7 +561,7 @@ jobs:
# shell: arch -arch arm64 bash {0}
# env:
# HF_HOME: /System/Volumes/Data/mnt/cache
# HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# HF_TOKEN: ${{ secrets.HF_TOKEN }}
# run: |
# ${CONDA_RUN} python -m pytest -n 1 -s -v --make-reports=tests_torch_mps \
# --report-log=tests_torch_mps.log \

View File

@@ -1,17 +0,0 @@
name: PR Style Bot
on:
issue_comment:
types: [created]
permissions:
contents: write
pull-requests: write
jobs:
style:
uses: huggingface/huggingface_hub/.github/workflows/style-bot-action.yml@main
with:
python_quality_dependencies: "[quality]"
secrets:
bot_token: ${{ secrets.GITHUB_TOKEN }}

View File

@@ -2,7 +2,8 @@ name: Fast tests for PRs
on:
pull_request:
branches: [main]
branches:
- main
paths:
- "src/diffusers/**.py"
- "benchmarks/**.py"
@@ -11,7 +12,6 @@ on:
- "tests/**.py"
- ".github/**.yml"
- "utils/**.py"
- "setup.py"
push:
branches:
- ci-*
@@ -64,7 +64,6 @@ jobs:
run: |
python utils/check_copies.py
python utils/check_dummies.py
python utils/check_support_list.py
make deps_table_check_updated
- name: Check if failure
if: ${{ failure() }}
@@ -121,8 +120,7 @@ jobs:
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git --no-deps
python -m uv pip install accelerate
- name: Environment
run: |
@@ -268,7 +266,6 @@ jobs:
# TODO (sayakpaul, DN6): revisit `--no-deps`
python -m pip install -U peft@git+https://github.com/huggingface/peft.git --no-deps
python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
python -m uv pip install -U tokenizers
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git --no-deps
- name: Environment

View File

@@ -1,296 +0,0 @@
name: Fast GPU Tests on PR
on:
pull_request:
branches: main
paths:
- "src/diffusers/models/modeling_utils.py"
- "src/diffusers/models/model_loading_utils.py"
- "src/diffusers/pipelines/pipeline_utils.py"
- "src/diffusers/pipeline_loading_utils.py"
- "src/diffusers/loaders/lora_base.py"
- "src/diffusers/loaders/lora_pipeline.py"
- "src/diffusers/loaders/peft.py"
- "tests/pipelines/test_pipelines_common.py"
- "tests/models/test_modeling_common.py"
workflow_dispatch:
concurrency:
group: ${{ github.workflow }}-${{ github.head_ref || github.run_id }}
cancel-in-progress: true
env:
DIFFUSERS_IS_CI: yes
OMP_NUM_THREADS: 8
MKL_NUM_THREADS: 8
HF_HUB_ENABLE_HF_TRANSFER: 1
PYTEST_TIMEOUT: 600
PIPELINE_USAGE_CUTOFF: 1000000000 # set high cutoff so that only always-test pipelines run
jobs:
check_code_quality:
runs-on: ubuntu-22.04
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check quality
run: make quality
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Quality check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make style && make quality'" >> $GITHUB_STEP_SUMMARY
check_repository_consistency:
needs: check_code_quality
runs-on: ubuntu-22.04
steps:
- uses: actions/checkout@v3
- name: Set up Python
uses: actions/setup-python@v4
with:
python-version: "3.8"
- name: Install dependencies
run: |
python -m pip install --upgrade pip
pip install .[quality]
- name: Check repo consistency
run: |
python utils/check_copies.py
python utils/check_dummies.py
python utils/check_support_list.py
make deps_table_check_updated
- name: Check if failure
if: ${{ failure() }}
run: |
echo "Repo consistency check failed. Please ensure the right dependency versions are installed with 'pip install -e .[quality]' and run 'make fix-copies'" >> $GITHUB_STEP_SUMMARY
setup_torch_cuda_pipeline_matrix:
needs: [check_code_quality, check_repository_consistency]
name: Setup Torch Pipelines CUDA Slow Tests Matrix
runs-on:
group: aws-general-8-plus
container:
image: diffusers/diffusers-pytorch-cpu
outputs:
pipeline_test_matrix: ${{ steps.fetch_pipeline_matrix.outputs.pipeline_test_matrix }}
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
- name: Environment
run: |
python utils/print_env.py
- name: Fetch Pipeline Matrix
id: fetch_pipeline_matrix
run: |
matrix=$(python utils/fetch_torch_cuda_pipeline_test_matrix.py)
echo $matrix
echo "pipeline_test_matrix=$matrix" >> $GITHUB_OUTPUT
- name: Pipeline Tests Artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: test-pipelines.json
path: reports
torch_pipelines_cuda_tests:
name: Torch Pipelines CUDA Tests
needs: setup_torch_cuda_pipeline_matrix
strategy:
fail-fast: false
max-parallel: 8
matrix:
module: ${{ fromJson(needs.setup_torch_cuda_pipeline_matrix.outputs.pipeline_test_matrix) }}
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
- name: Environment
run: |
python utils/print_env.py
- name: Extract tests
id: extract_tests
run: |
pattern=$(python utils/extract_tests_from_mixin.py --type pipeline)
echo "$pattern" > /tmp/test_pattern.txt
echo "pattern_file=/tmp/test_pattern.txt" >> $GITHUB_OUTPUT
- name: PyTorch CUDA checkpoint tests on Ubuntu
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
if [ "${{ matrix.module }}" = "ip_adapters" ]; then
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
tests/pipelines/${{ matrix.module }}
else
pattern=$(cat ${{ steps.extract_tests.outputs.pattern_file }})
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx and $pattern" \
--make-reports=tests_pipeline_${{ matrix.module }}_cuda \
tests/pipelines/${{ matrix.module }}
fi
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_pipeline_${{ matrix.module }}_cuda_stats.txt
cat reports/tests_pipeline_${{ matrix.module }}_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: pipeline_${{ matrix.module }}_test_reports
path: reports
torch_cuda_tests:
name: Torch CUDA Tests
needs: [check_code_quality, check_repository_consistency]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --shm-size "16gb" --ipc host --gpus 0
defaults:
run:
shell: bash
strategy:
fail-fast: false
max-parallel: 2
matrix:
module: [models, schedulers, lora, others]
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install -e [quality,test]
python -m uv pip install peft@git+https://github.com/huggingface/peft.git
pip uninstall accelerate -y && python -m uv pip install -U accelerate@git+https://github.com/huggingface/accelerate.git
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
- name: Environment
run: |
python utils/print_env.py
- name: Extract tests
id: extract_tests
run: |
pattern=$(python utils/extract_tests_from_mixin.py --type ${{ matrix.module }})
echo "$pattern" > /tmp/test_pattern.txt
echo "pattern_file=/tmp/test_pattern.txt" >> $GITHUB_OUTPUT
- name: Run PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
pattern=$(cat ${{ steps.extract_tests.outputs.pattern_file }})
if [ -z "$pattern" ]; then
python -m pytest -n 1 -sv --max-worker-restart=0 --dist=loadfile -k "not Flax and not Onnx" tests/${{ matrix.module }} \
--make-reports=tests_torch_cuda_${{ matrix.module }}
else
python -m pytest -n 1 -sv --max-worker-restart=0 --dist=loadfile -k "not Flax and not Onnx and $pattern" tests/${{ matrix.module }} \
--make-reports=tests_torch_cuda_${{ matrix.module }}
fi
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/tests_torch_cuda_${{ matrix.module }}_stats.txt
cat reports/tests_torch_cuda_${{ matrix.module }}_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: torch_cuda_test_reports_${{ matrix.module }}
path: reports
run_examples_tests:
name: Examples PyTorch CUDA tests on Ubuntu
needs: [check_code_quality, check_repository_consistency]
runs-on:
group: aws-g4dn-2xlarge
container:
image: diffusers/diffusers-pytorch-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
with:
fetch-depth: 2
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
pip uninstall transformers -y && python -m uv pip install -U transformers@git+https://github.com/huggingface/transformers.git --no-deps
python -m uv pip install -e [quality,test,training]
- name: Environment
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install timm
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v --make-reports=examples_torch_cuda examples/
- name: Failure short reports
if: ${{ failure() }}
run: |
cat reports/examples_torch_cuda_stats.txt
cat reports/examples_torch_cuda_failures_short.txt
- name: Test suite reports artifacts
if: ${{ always() }}
uses: actions/upload-artifact@v4
with:
name: examples_test_reports
path: reports

View File

@@ -83,7 +83,7 @@ jobs:
python utils/print_env.py
- name: PyTorch CUDA checkpoint tests on Ubuntu
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
@@ -137,7 +137,7 @@ jobs:
- name: Run PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
@@ -187,7 +187,7 @@ jobs:
- name: Run Flax TPU tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
@@ -235,7 +235,7 @@ jobs:
- name: Run ONNXRuntime CUDA tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
@@ -283,7 +283,7 @@ jobs:
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
RUN_COMPILE: yes
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "compile" --make-reports=tests_torch_compile_cuda tests/
@@ -326,7 +326,7 @@ jobs:
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "xformers" --make-reports=tests_torch_xformers_cuda tests/
- name: Failure short reports
@@ -349,6 +349,7 @@ jobs:
container:
image: diffusers/diffusers-pytorch-cuda
options: --gpus 0 --shm-size "16gb" --ipc host
steps:
- name: Checkout diffusers
uses: actions/checkout@v3
@@ -358,6 +359,7 @@ jobs:
- name: NVIDIA-SMI
run: |
nvidia-smi
- name: Install dependencies
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
@@ -370,7 +372,7 @@ jobs:
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install timm

View File

@@ -81,7 +81,7 @@ jobs:
python utils/print_env.py
- name: Slow PyTorch CUDA checkpoint tests on Ubuntu
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
@@ -135,7 +135,7 @@ jobs:
- name: Run PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
@@ -186,14 +186,14 @@ jobs:
- name: Run PyTorch CUDA tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
# https://pytorch.org/docs/stable/notes/randomness.html#avoiding-nondeterministic-algorithms
CUBLAS_WORKSPACE_CONFIG: :16:8
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "not Flax and not Onnx" \
--make-reports=tests_torch_minimum_cuda \
tests/models/test_modeling_common.py \
tests/models/test_modelling_common.py \
tests/pipelines/test_pipelines_common.py \
tests/pipelines/test_pipeline_utils.py \
tests/pipelines/test_pipelines.py \
@@ -241,7 +241,7 @@ jobs:
- name: Run slow Flax TPU tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 0 \
-s -v -k "Flax" \
@@ -289,7 +289,7 @@ jobs:
- name: Run slow ONNXRuntime CUDA tests
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile \
-s -v -k "Onnx" \
@@ -335,9 +335,9 @@ jobs:
- name: Environment
run: |
python utils/print_env.py
- name: Run torch compile tests on GPU
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
RUN_COMPILE: yes
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "compile" --make-reports=tests_torch_compile_cuda tests/
@@ -380,7 +380,7 @@ jobs:
python utils/print_env.py
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m pytest -n 1 --max-worker-restart=0 --dist=loadfile -s -v -k "xformers" --make-reports=tests_torch_xformers_cuda tests/
- name: Failure short reports
@@ -426,7 +426,7 @@ jobs:
- name: Run example tests on GPU
env:
HF_TOKEN: ${{ secrets.DIFFUSERS_HF_HUB_READ_TOKEN }}
HF_TOKEN: ${{ secrets.HF_TOKEN }}
run: |
python -m venv /opt/venv && export PATH="/opt/venv/bin:$PATH"
python -m uv pip install timm

View File

@@ -7,8 +7,8 @@ on:
default: 'diffusers/diffusers-pytorch-cuda'
description: 'Name of the Docker image'
required: true
pr_number:
description: 'PR number to test on'
branch:
description: 'PR Branch to test on'
required: true
test:
description: 'Tests to run (e.g.: `tests/models`).'
@@ -43,8 +43,8 @@ jobs:
exit 1
fi
if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines|lora) ]]; then
echo "Error: The input string must contain either 'models', 'pipelines', or 'lora' after 'tests/'."
if [[ ! "$PY_TEST" =~ ^tests/(models|pipelines) ]]; then
echo "Error: The input string must contain either 'models' or 'pipelines' after 'tests/'."
exit 1
fi
@@ -53,13 +53,13 @@ jobs:
exit 1
fi
echo "$PY_TEST"
shell: bash -e {0}
- name: Checkout PR branch
uses: actions/checkout@v4
with:
ref: refs/pull/${{ inputs.pr_number }}/head
ref: ${{ github.event.inputs.branch }}
repository: ${{ github.event.pull_request.head.repo.full_name }}
- name: Install pytest
run: |

View File

@@ -13,6 +13,3 @@ jobs:
fetch-depth: 0
- name: Secret Scanning
uses: trufflesecurity/trufflehog@main
with:
extra_args: --results=verified,unknown

View File

@@ -28,9 +28,9 @@ ENV PATH="/opt/venv/bin:$PATH"
# pre-install the heavy dependencies (these can later be overridden by the deps from setup.py)
RUN python3 -m pip install --no-cache-dir --upgrade pip uv==0.1.11 && \
python3 -m uv pip install --no-cache-dir \
torch \
torchvision \
torchaudio\
torch==2.1.2 \
torchvision==0.16.2 \
torchaudio==2.1.2 \
onnxruntime \
--extra-index-url https://download.pytorch.org/whl/cpu && \
python3 -m uv pip install --no-cache-dir \

View File

@@ -17,8 +17,12 @@
title: AutoPipeline
- local: tutorials/basic_training
title: Train a diffusion model
- local: tutorials/using_peft_for_inference
title: Load LoRAs for inference
- local: tutorials/fast_diffusion
title: Accelerate inference of text-to-image diffusion models
- local: tutorials/inference_with_big_models
title: Working with big models
title: Tutorials
- sections:
- local: using-diffusers/loading
@@ -29,24 +33,11 @@
title: Load schedulers and models
- local: using-diffusers/other-formats
title: Model files and layouts
- local: using-diffusers/loading_adapters
title: Load adapters
- local: using-diffusers/push_to_hub
title: Push files to the Hub
title: Load pipelines and adapters
- sections:
- local: tutorials/using_peft_for_inference
title: LoRA
- local: using-diffusers/ip_adapter
title: IP-Adapter
- local: using-diffusers/controlnet
title: ControlNet
- local: using-diffusers/t2i_adapter
title: T2I-Adapter
- local: using-diffusers/dreambooth
title: DreamBooth
- local: using-diffusers/textual_inversion_inference
title: Textual inversion
title: Adapters
isExpanded: false
- sections:
- local: using-diffusers/unconditional_image_generation
title: Unconditional image generation
@@ -68,6 +59,8 @@
title: Create a server
- local: training/distributed_inference
title: Distributed inference
- local: using-diffusers/merge_loras
title: Merge LoRAs
- local: using-diffusers/scheduler_features
title: Scheduler features
- local: using-diffusers/callback
@@ -83,33 +76,27 @@
- local: advanced_inference/outpaint
title: Outpainting
title: Advanced inference
- sections:
- local: hybrid_inference/overview
title: Overview
- local: hybrid_inference/vae_decode
title: VAE Decode
- local: hybrid_inference/vae_encode
title: VAE Encode
- local: hybrid_inference/api_reference
title: API Reference
title: Hybrid Inference
- sections:
- local: using-diffusers/cogvideox
title: CogVideoX
- local: using-diffusers/consisid
title: ConsisID
- local: using-diffusers/sdxl
title: Stable Diffusion XL
- local: using-diffusers/sdxl_turbo
title: SDXL Turbo
- local: using-diffusers/kandinsky
title: Kandinsky
- local: using-diffusers/omnigen
title: OmniGen
- local: using-diffusers/ip_adapter
title: IP-Adapter
- local: using-diffusers/pag
title: PAG
- local: using-diffusers/controlnet
title: ControlNet
- local: using-diffusers/t2i_adapter
title: T2I-Adapter
- local: using-diffusers/inference_with_lcm
title: Latent Consistency Model
- local: using-diffusers/textual_inversion_inference
title: Textual inversion
- local: using-diffusers/shap-e
title: Shap-E
- local: using-diffusers/diffedit
@@ -174,12 +161,10 @@
title: gguf
- local: quantization/torchao
title: torchao
- local: quantization/quanto
title: quanto
title: Quantization Methods
- sections:
- local: optimization/fp16
title: Accelerate inference
title: Speed up inference
- local: optimization/memory
title: Reduce memory usage
- local: optimization/torch2.0
@@ -194,8 +179,6 @@
title: TGATE
- local: optimization/xdit
title: xDiT
- local: optimization/para_attn
title: ParaAttention
- sections:
- local: using-diffusers/stable_diffusion_jax_how_to
title: JAX/Flax
@@ -264,23 +247,19 @@
sections:
- local: api/models/overview
title: Overview
- local: api/models/auto_model
title: AutoModel
- sections:
- local: api/models/controlnet
title: ControlNetModel
- local: api/models/controlnet_union
title: ControlNetUnionModel
- local: api/models/controlnet_flux
title: FluxControlNetModel
- local: api/models/controlnet_hunyuandit
title: HunyuanDiT2DControlNetModel
- local: api/models/controlnet_sana
title: SanaControlNetModel
- local: api/models/controlnet_sd3
title: SD3ControlNetModel
- local: api/models/controlnet_sparsectrl
title: SparseControlNetModel
- local: api/models/controlnet_union
title: ControlNetUnionModel
title: ControlNets
- sections:
- local: api/models/allegro_transformer3d
@@ -291,62 +270,46 @@
title: CogVideoXTransformer3DModel
- local: api/models/cogview3plus_transformer2d
title: CogView3PlusTransformer2DModel
- local: api/models/cogview4_transformer2d
title: CogView4Transformer2DModel
- local: api/models/consisid_transformer3d
title: ConsisIDTransformer3DModel
- local: api/models/cosmos_transformer3d
title: CosmosTransformer3DModel
- local: api/models/dit_transformer2d
title: DiTTransformer2DModel
- local: api/models/easyanimate_transformer3d
title: EasyAnimateTransformer3DModel
- local: api/models/flux_transformer
title: FluxTransformer2DModel
- local: api/models/hidream_image_transformer
title: HiDreamImageTransformer2DModel
- local: api/models/hunyuan_transformer2d
title: HunyuanDiT2DModel
- local: api/models/hunyuan_video_transformer_3d
title: HunyuanVideoTransformer3DModel
- local: api/models/latte_transformer3d
title: LatteTransformer3DModel
- local: api/models/ltx_video_transformer3d
title: LTXVideoTransformer3DModel
- local: api/models/lumina2_transformer2d
title: Lumina2Transformer2DModel
- local: api/models/lumina_nextdit2d
title: LuminaNextDiT2DModel
- local: api/models/ltx_video_transformer3d
title: LTXVideoTransformer3DModel
- local: api/models/mochi_transformer3d
title: MochiTransformer3DModel
- local: api/models/omnigen_transformer
title: OmniGenTransformer2DModel
- local: api/models/pixart_transformer2d
title: PixArtTransformer2DModel
- local: api/models/prior_transformer
title: PriorTransformer
- local: api/models/sana_transformer2d
title: SanaTransformer2DModel
- local: api/models/sd3_transformer2d
title: SD3Transformer2DModel
- local: api/models/sana_transformer2d
title: SanaTransformer2DModel
- local: api/models/stable_audio_transformer
title: StableAudioDiTModel
- local: api/models/transformer2d
title: Transformer2DModel
- local: api/models/transformer_temporal
title: TransformerTemporalModel
- local: api/models/wan_transformer_3d
title: WanTransformer3DModel
title: Transformers
- sections:
- local: api/models/stable_cascade_unet
title: StableCascadeUNet
- local: api/models/unet
title: UNet1DModel
- local: api/models/unet2d-cond
title: UNet2DConditionModel
- local: api/models/unet2d
title: UNet2DModel
- local: api/models/unet2d-cond
title: UNet2DConditionModel
- local: api/models/unet3d-cond
title: UNet3DConditionModel
- local: api/models/unet-motion
@@ -355,28 +318,22 @@
title: UViT2DModel
title: UNets
- sections:
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_dc
title: AutoencoderDC
- local: api/models/autoencoderkl
title: AutoencoderKL
- local: api/models/autoencoderkl_allegro
title: AutoencoderKLAllegro
- local: api/models/autoencoderkl_cogvideox
title: AutoencoderKLCogVideoX
- local: api/models/autoencoderkl_cosmos
title: AutoencoderKLCosmos
- local: api/models/autoencoder_kl_hunyuan_video
title: AutoencoderKLHunyuanVideo
- local: api/models/autoencoderkl_ltx_video
title: AutoencoderKLLTXVideo
- local: api/models/autoencoderkl_magvit
title: AutoencoderKLMagvit
- local: api/models/autoencoderkl_mochi
title: AutoencoderKLMochi
- local: api/models/autoencoder_kl_wan
title: AutoencoderKLWan
- local: api/models/asymmetricautoencoderkl
title: AsymmetricAutoencoderKL
- local: api/models/autoencoder_dc
title: AutoencoderDC
- local: api/models/consistency_decoder_vae
title: ConsistencyDecoderVAE
- local: api/models/autoencoder_oobleck
@@ -413,10 +370,6 @@
title: CogVideoX
- local: api/pipelines/cogview3
title: CogView3
- local: api/pipelines/cogview4
title: CogView4
- local: api/pipelines/consisid
title: ConsisID
- local: api/pipelines/consistency_models
title: Consistency Models
- local: api/pipelines/controlnet
@@ -429,16 +382,12 @@
title: ControlNet with Stable Diffusion 3
- local: api/pipelines/controlnet_sdxl
title: ControlNet with Stable Diffusion XL
- local: api/pipelines/controlnet_sana
title: ControlNet-Sana
- local: api/pipelines/controlnetxs
title: ControlNet-XS
- local: api/pipelines/controlnetxs_sdxl
title: ControlNet-XS with Stable Diffusion XL
- local: api/pipelines/controlnet_union
title: ControlNetUnion
- local: api/pipelines/cosmos
title: Cosmos
- local: api/pipelines/dance_diffusion
title: Dance Diffusion
- local: api/pipelines/ddim
@@ -451,16 +400,10 @@
title: DiffEdit
- local: api/pipelines/dit
title: DiT
- local: api/pipelines/easyanimate
title: EasyAnimate
- local: api/pipelines/flux
title: Flux
- local: api/pipelines/control_flux_inpaint
title: FluxControlInpaint
- local: api/pipelines/framepack
title: Framepack
- local: api/pipelines/hidream
title: HiDream-I1
- local: api/pipelines/hunyuandit
title: Hunyuan-DiT
- local: api/pipelines/hunyuan_video
@@ -487,8 +430,6 @@
title: LEDITS++
- local: api/pipelines/ltx_video
title: LTXVideo
- local: api/pipelines/lumina2
title: Lumina 2.0
- local: api/pipelines/lumina
title: Lumina-T2X
- local: api/pipelines/marigold
@@ -499,8 +440,6 @@
title: MultiDiffusion
- local: api/pipelines/musicldm
title: MusicLDM
- local: api/pipelines/omnigen
title: OmniGen
- local: api/pipelines/pag
title: PAG
- local: api/pipelines/paint_by_example
@@ -513,8 +452,6 @@
title: PixArt-Σ
- local: api/pipelines/sana
title: Sana
- local: api/pipelines/sana_sprint
title: Sana Sprint
- local: api/pipelines/self_attention_guidance
title: Self-Attention Guidance
- local: api/pipelines/semantic_stable_diffusion
@@ -528,40 +465,40 @@
- sections:
- local: api/pipelines/stable_diffusion/overview
title: Overview
- local: api/pipelines/stable_diffusion/depth2img
title: Depth-to-image
- local: api/pipelines/stable_diffusion/gligen
title: GLIGEN (Grounded Language-to-Image Generation)
- local: api/pipelines/stable_diffusion/image_variation
title: Image variation
- local: api/pipelines/stable_diffusion/text2img
title: Text-to-image
- local: api/pipelines/stable_diffusion/img2img
title: Image-to-image
- local: api/pipelines/stable_diffusion/svd
title: Image-to-video
- local: api/pipelines/stable_diffusion/inpaint
title: Inpainting
- local: api/pipelines/stable_diffusion/k_diffusion
title: K-Diffusion
- local: api/pipelines/stable_diffusion/latent_upscale
title: Latent upscaler
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D Upscaler
- local: api/pipelines/stable_diffusion/depth2img
title: Depth-to-image
- local: api/pipelines/stable_diffusion/image_variation
title: Image variation
- local: api/pipelines/stable_diffusion/stable_diffusion_safe
title: Safe Stable Diffusion
- local: api/pipelines/stable_diffusion/sdxl_turbo
title: SDXL Turbo
- local: api/pipelines/stable_diffusion/stable_diffusion_2
title: Stable Diffusion 2
- local: api/pipelines/stable_diffusion/stable_diffusion_3
title: Stable Diffusion 3
- local: api/pipelines/stable_diffusion/stable_diffusion_xl
title: Stable Diffusion XL
- local: api/pipelines/stable_diffusion/sdxl_turbo
title: SDXL Turbo
- local: api/pipelines/stable_diffusion/latent_upscale
title: Latent upscaler
- local: api/pipelines/stable_diffusion/upscale
title: Super-resolution
- local: api/pipelines/stable_diffusion/k_diffusion
title: K-Diffusion
- local: api/pipelines/stable_diffusion/ldm3d_diffusion
title: LDM3D Text-to-(RGB, Depth), Text-to-(RGB-pano, Depth-pano), LDM3D Upscaler
- local: api/pipelines/stable_diffusion/adapter
title: T2I-Adapter
- local: api/pipelines/stable_diffusion/text2img
title: Text-to-image
- local: api/pipelines/stable_diffusion/gligen
title: GLIGEN (Grounded Language-to-Image Generation)
title: Stable Diffusion
- local: api/pipelines/stable_unclip
title: Stable unCLIP
@@ -575,10 +512,6 @@
title: UniDiffuser
- local: api/pipelines/value_guided_sampling
title: Value-guided sampling
- local: api/pipelines/visualcloze
title: VisualCloze
- local: api/pipelines/wan
title: Wan
- local: api/pipelines/wuerstchen
title: Wuerstchen
title: Pipelines
@@ -588,10 +521,6 @@
title: Overview
- local: api/schedulers/cm_stochastic_iterative
title: CMStochasticIterativeScheduler
- local: api/schedulers/ddim_cogvideox
title: CogVideoXDDIMScheduler
- local: api/schedulers/multistep_dpm_solver_cogvideox
title: CogVideoXDPMScheduler
- local: api/schedulers/consistency_decoder
title: ConsistencyDecoderScheduler
- local: api/schedulers/cosine_dpm
@@ -661,8 +590,6 @@
title: Attention Processor
- local: api/activations
title: Custom activation functions
- local: api/cache
title: Caching methods
- local: api/normalization
title: Custom normalization layers
- local: api/utilities

View File

@@ -25,16 +25,3 @@ Customized activation functions for supporting various models in 🤗 Diffusers.
## ApproximateGELU
[[autodoc]] models.activations.ApproximateGELU
## SwiGLU
[[autodoc]] models.activations.SwiGLU
## FP32SiLU
[[autodoc]] models.activations.FP32SiLU
## LinearActivation
[[autodoc]] models.activations.LinearActivation

View File

@@ -147,20 +147,3 @@ An attention processor is a class for applying different types of attention mech
## XLAFlashAttnProcessor2_0
[[autodoc]] models.attention_processor.XLAFlashAttnProcessor2_0
## XFormersJointAttnProcessor
[[autodoc]] models.attention_processor.XFormersJointAttnProcessor
## IPAdapterXFormersAttnProcessor
[[autodoc]] models.attention_processor.IPAdapterXFormersAttnProcessor
## FluxIPAdapterJointAttnProcessor2_0
[[autodoc]] models.attention_processor.FluxIPAdapterJointAttnProcessor2_0
## XLAFluxFlashAttnProcessor2_0
[[autodoc]] models.attention_processor.XLAFluxFlashAttnProcessor2_0

View File

@@ -1,82 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# Caching methods
## Pyramid Attention Broadcast
[Pyramid Attention Broadcast](https://huggingface.co/papers/2408.12588) from Xuanlei Zhao, Xiaolong Jin, Kai Wang, Yang You.
Pyramid Attention Broadcast (PAB) is a method that speeds up inference in diffusion models by systematically skipping attention computations between successive inference steps and reusing cached attention states. The attention states are not very different between successive inference steps. The most prominent difference is in the spatial attention blocks, not as much in the temporal attention blocks, and finally the least in the cross attention blocks. Therefore, many cross attention computation blocks can be skipped, followed by the temporal and spatial attention blocks. By combining other techniques like sequence parallelism and classifier-free guidance parallelism, PAB achieves near real-time video generation.
Enable PAB with [`~PyramidAttentionBroadcastConfig`] on any pipeline. For some benchmarks, refer to [this](https://github.com/huggingface/diffusers/pull/9562) pull request.
```python
import torch
from diffusers import CogVideoXPipeline, PyramidAttentionBroadcastConfig
pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
pipe.to("cuda")
# Increasing the value of `spatial_attention_timestep_skip_range[0]` or decreasing the value of
# `spatial_attention_timestep_skip_range[1]` will decrease the interval in which pyramid attention
# broadcast is active, leader to slower inference speeds. However, large intervals can lead to
# poorer quality of generated videos.
config = PyramidAttentionBroadcastConfig(
spatial_attention_block_skip_range=2,
spatial_attention_timestep_skip_range=(100, 800),
current_timestep_callback=lambda: pipe.current_timestep,
)
pipe.transformer.enable_cache(config)
```
## Faster Cache
[FasterCache](https://huggingface.co/papers/2410.19355) from Zhengyao Lv, Chenyang Si, Junhao Song, Zhenyu Yang, Yu Qiao, Ziwei Liu, Kwan-Yee K. Wong.
FasterCache is a method that speeds up inference in diffusion transformers by:
- Reusing attention states between successive inference steps, due to high similarity between them
- Skipping unconditional branch prediction used in classifier-free guidance by revealing redundancies between unconditional and conditional branch outputs for the same timestep, and therefore approximating the unconditional branch output using the conditional branch output
```python
import torch
from diffusers import CogVideoXPipeline, FasterCacheConfig
pipe = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
pipe.to("cuda")
config = FasterCacheConfig(
spatial_attention_block_skip_range=2,
spatial_attention_timestep_skip_range=(-1, 681),
current_timestep_callback=lambda: pipe.current_timestep,
attention_weight_callback=lambda _: 0.3,
unconditional_batch_skip_range=5,
unconditional_batch_timestep_skip_range=(-1, 781),
tensor_format="BFCHW",
)
pipe.transformer.enable_cache(config)
```
### CacheMixin
[[autodoc]] CacheMixin
### PyramidAttentionBroadcastConfig
[[autodoc]] PyramidAttentionBroadcastConfig
[[autodoc]] apply_pyramid_attention_broadcast
### FasterCacheConfig
[[autodoc]] FasterCacheConfig
[[autodoc]] apply_faster_cache

View File

@@ -20,15 +20,7 @@ LoRA is a fast and lightweight training method that inserts and trains a signifi
- [`FluxLoraLoaderMixin`] provides similar functions for [Flux](https://huggingface.co/docs/diffusers/main/en/api/pipelines/flux).
- [`CogVideoXLoraLoaderMixin`] provides similar functions for [CogVideoX](https://huggingface.co/docs/diffusers/main/en/api/pipelines/cogvideox).
- [`Mochi1LoraLoaderMixin`] provides similar functions for [Mochi](https://huggingface.co/docs/diffusers/main/en/api/pipelines/mochi).
- [`AuraFlowLoraLoaderMixin`] provides similar functions for [AuraFlow](https://huggingface.co/fal/AuraFlow).
- [`LTXVideoLoraLoaderMixin`] provides similar functions for [LTX-Video](https://huggingface.co/docs/diffusers/main/en/api/pipelines/ltx_video).
- [`SanaLoraLoaderMixin`] provides similar functions for [Sana](https://huggingface.co/docs/diffusers/main/en/api/pipelines/sana).
- [`HunyuanVideoLoraLoaderMixin`] provides similar functions for [HunyuanVideo](https://huggingface.co/docs/diffusers/main/en/api/pipelines/hunyuan_video).
- [`Lumina2LoraLoaderMixin`] provides similar functions for [Lumina2](https://huggingface.co/docs/diffusers/main/en/api/pipelines/lumina2).
- [`WanLoraLoaderMixin`] provides similar functions for [Wan](https://huggingface.co/docs/diffusers/main/en/api/pipelines/wan).
- [`CogView4LoraLoaderMixin`] provides similar functions for [CogView4](https://huggingface.co/docs/diffusers/main/en/api/pipelines/cogview4).
- [`AmusedLoraLoaderMixin`] is for the [`AmusedPipeline`].
- [`HiDreamImageLoraLoaderMixin`] provides similar functions for [HiDream Image](https://huggingface.co/docs/diffusers/main/en/api/pipelines/hidream)
- [`LoraBaseMixin`] provides a base class with several utility methods to fuse, unfuse, unload, LoRAs and more.
<Tip>
@@ -60,42 +52,11 @@ To learn more about how to load LoRA weights, see the [LoRA](../../using-diffuse
## Mochi1LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.Mochi1LoraLoaderMixin
## AuraFlowLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.AuraFlowLoraLoaderMixin
## LTXVideoLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.LTXVideoLoraLoaderMixin
## SanaLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.SanaLoraLoaderMixin
## HunyuanVideoLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.HunyuanVideoLoraLoaderMixin
## Lumina2LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.Lumina2LoraLoaderMixin
## CogView4LoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.CogView4LoraLoaderMixin
## WanLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.WanLoraLoaderMixin
## AmusedLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.AmusedLoraLoaderMixin
## HiDreamImageLoraLoaderMixin
[[autodoc]] loaders.lora_pipeline.HiDreamImageLoraLoaderMixin
## LoraBaseMixin
[[autodoc]] loaders.lora_base.LoraBaseMixin

View File

@@ -1,29 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# AutoModel
The `AutoModel` is designed to make it easy to load a checkpoint without needing to know the specific model class. `AutoModel` automatically retrieves the correct model class from the checkpoint `config.json` file.
```python
from diffusers import AutoModel, AutoPipelineForText2Image
unet = AutoModel.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", subfolder="unet")
pipe = AutoPipelineForText2Image.from_pretrained("stable-diffusion-v1-5/stable-diffusion-v1-5", unet=unet)
```
## AutoModel
[[autodoc]] AutoModel
- all
- from_pretrained

View File

@@ -1,32 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLWan
The 3D variational autoencoder (VAE) model with KL loss used in [Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLWan
vae = AutoencoderKLWan.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", subfolder="vae", torch_dtype=torch.float32)
```
## AutoencoderKLWan
[[autodoc]] AutoencoderKLWan
- decode
- all
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

View File

@@ -18,7 +18,7 @@ The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLAllegro
vae = AutoencoderKLAllegro.from_pretrained("rhymes-ai/Allegro", subfolder="vae", torch_dtype=torch.float32).to("cuda")
vae = AutoencoderKLCogVideoX.from_pretrained("rhymes-ai/Allegro", subfolder="vae", torch_dtype=torch.float32).to("cuda")
```
## AutoencoderKLAllegro

View File

@@ -1,40 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLCosmos
[Cosmos Tokenizers](https://github.com/NVIDIA/Cosmos-Tokenizer).
Supported models:
- [nvidia/Cosmos-1.0-Tokenizer-CV8x8x8](https://huggingface.co/nvidia/Cosmos-1.0-Tokenizer-CV8x8x8)
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLCosmos
vae = AutoencoderKLCosmos.from_pretrained("nvidia/Cosmos-1.0-Tokenizer-CV8x8x8", subfolder="vae")
```
## AutoencoderKLCosmos
[[autodoc]] AutoencoderKLCosmos
- decode
- encode
- all
## AutoencoderKLOutput
[[autodoc]] models.autoencoders.autoencoder_kl.AutoencoderKLOutput
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

View File

@@ -1,37 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# AutoencoderKLMagvit
The 3D variational autoencoder (VAE) model with KL loss used in [EasyAnimate](https://github.com/aigc-apps/EasyAnimate) was introduced by Alibaba PAI.
The model can be loaded with the following code snippet.
```python
from diffusers import AutoencoderKLMagvit
vae = AutoencoderKLMagvit.from_pretrained("alibaba-pai/EasyAnimateV5.1-12b-zh", subfolder="vae", torch_dtype=torch.float16).to("cuda")
```
## AutoencoderKLMagvit
[[autodoc]] AutoencoderKLMagvit
- decode
- encode
- all
## AutoencoderKLOutput
[[autodoc]] models.autoencoders.autoencoder_kl.AutoencoderKLOutput
## DecoderOutput
[[autodoc]] models.autoencoders.vae.DecoderOutput

View File

@@ -1,30 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# CogView4Transformer2DModel
A Diffusion Transformer model for 2D data from [CogView4]()
The model can be loaded with the following code snippet.
```python
from diffusers import CogView4Transformer2DModel
transformer = CogView4Transformer2DModel.from_pretrained("THUDM/CogView4-6B", subfolder="transformer", torch_dtype=torch.bfloat16).to("cuda")
```
## CogView4Transformer2DModel
[[autodoc]] CogView4Transformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,30 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# ConsisIDTransformer3DModel
A Diffusion Transformer model for 3D data from [ConsisID](https://github.com/PKU-YuanGroup/ConsisID) was introduced in [Identity-Preserving Text-to-Video Generation by Frequency Decomposition](https://arxiv.org/pdf/2411.17440) by Peking University & University of Rochester & etc.
The model can be loaded with the following code snippet.
```python
from diffusers import ConsisIDTransformer3DModel
transformer = ConsisIDTransformer3DModel.from_pretrained("BestWishYsh/ConsisID-preview", subfolder="transformer", torch_dtype=torch.bfloat16).to("cuda")
```
## ConsisIDTransformer3DModel
[[autodoc]] ConsisIDTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,29 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# SanaControlNetModel
The ControlNet model was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, Maneesh Agrawala. It provides a greater degree of control over text-to-image generation by conditioning the model on additional inputs such as edge maps, depth maps, segmentation maps, and keypoints for pose detection.
The abstract from the paper is:
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
This model was contributed by [ishan24](https://huggingface.co/ishan24). ❤️
The original codebase can be found at [NVlabs/Sana](https://github.com/NVlabs/Sana), and you can find official ControlNet checkpoints on [Efficient-Large-Model's](https://huggingface.co/Efficient-Large-Model) Hub profile.
## SanaControlNetModel
[[autodoc]] SanaControlNetModel
## SanaControlNetOutput
[[autodoc]] models.controlnets.controlnet_sana.SanaControlNetOutput

View File

@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# CosmosTransformer3DModel
A Diffusion Transformer model for 3D video-like data was introduced in [Cosmos World Foundation Model Platform for Physical AI](https://huggingface.co/papers/2501.03575) by NVIDIA.
The model can be loaded with the following code snippet.
```python
from diffusers import CosmosTransformer3DModel
transformer = CosmosTransformer3DModel.from_pretrained("nvidia/Cosmos-1.0-Diffusion-7B-Text2World", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## CosmosTransformer3DModel
[[autodoc]] CosmosTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,30 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# EasyAnimateTransformer3DModel
A Diffusion Transformer model for 3D data from [EasyAnimate](https://github.com/aigc-apps/EasyAnimate) was introduced by Alibaba PAI.
The model can be loaded with the following code snippet.
```python
from diffusers import EasyAnimateTransformer3DModel
transformer = EasyAnimateTransformer3DModel.from_pretrained("alibaba-pai/EasyAnimateV5.1-12b-zh", subfolder="transformer", torch_dtype=torch.float16).to("cuda")
```
## EasyAnimateTransformer3DModel
[[autodoc]] EasyAnimateTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,46 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# HiDreamImageTransformer2DModel
A Transformer model for image-like data from [HiDream-I1](https://huggingface.co/HiDream-ai).
The model can be loaded with the following code snippet.
```python
from diffusers import HiDreamImageTransformer2DModel
transformer = HiDreamImageTransformer2DModel.from_pretrained("HiDream-ai/HiDream-I1-Full", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## Loading GGUF quantized checkpoints for HiDream-I1
GGUF checkpoints for the `HiDreamImageTransformer2DModel` can be loaded using `~FromOriginalModelMixin.from_single_file`
```python
import torch
from diffusers import GGUFQuantizationConfig, HiDreamImageTransformer2DModel
ckpt_path = "https://huggingface.co/city96/HiDream-I1-Dev-gguf/blob/main/hidream-i1-dev-Q2_K.gguf"
transformer = HiDreamImageTransformer2DModel.from_single_file(
ckpt_path,
quantization_config=GGUFQuantizationConfig(compute_dtype=torch.bfloat16),
torch_dtype=torch.bfloat16
)
```
## HiDreamImageTransformer2DModel
[[autodoc]] HiDreamImageTransformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# Lumina2Transformer2DModel
A Diffusion Transformer model for 3D video-like data was introduced in [Lumina Image 2.0](https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0) by Alpha-VLLM.
The model can be loaded with the following code snippet.
```python
from diffusers import Lumina2Transformer2DModel
transformer = Lumina2Transformer2DModel.from_pretrained("Alpha-VLLM/Lumina-Image-2.0", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## Lumina2Transformer2DModel
[[autodoc]] Lumina2Transformer2DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -1,30 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# OmniGenTransformer2DModel
A Transformer model that accepts multimodal instructions to generate images for [OmniGen](https://github.com/VectorSpaceLab/OmniGen/).
The abstract from the paper is:
*The emergence of Large Language Models (LLMs) has unified language generation tasks and revolutionized human-machine interaction. However, in the realm of image generation, a unified model capable of handling various tasks within a single framework remains largely unexplored. In this work, we introduce OmniGen, a new diffusion model for unified image generation. OmniGen is characterized by the following features: 1) Unification: OmniGen not only demonstrates text-to-image generation capabilities but also inherently supports various downstream tasks, such as image editing, subject-driven generation, and visual conditional generation. 2) Simplicity: The architecture of OmniGen is highly simplified, eliminating the need for additional plugins. Moreover, compared to existing diffusion models, it is more user-friendly and can complete complex tasks end-to-end through instructions without the need for extra intermediate steps, greatly simplifying the image generation workflow. 3) Knowledge Transfer: Benefit from learning in a unified format, OmniGen effectively transfers knowledge across different tasks, manages unseen tasks and domains, and exhibits novel capabilities. We also explore the models reasoning capabilities and potential applications of the chain-of-thought mechanism. This work represents the first attempt at a general-purpose image generation model, and we will release our resources at https://github.com/VectorSpaceLab/OmniGen to foster future advancements.*
```python
import torch
from diffusers import OmniGenTransformer2DModel
transformer = OmniGenTransformer2DModel.from_pretrained("Shitao/OmniGen-v1-diffusers", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## OmniGenTransformer2DModel
[[autodoc]] OmniGenTransformer2DModel

View File

@@ -1,30 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License. -->
# WanTransformer3DModel
A Diffusion Transformer model for 3D video-like data was introduced in [Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
The model can be loaded with the following code snippet.
```python
from diffusers import WanTransformer3DModel
transformer = WanTransformer3DModel.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", subfolder="transformer", torch_dtype=torch.bfloat16)
```
## WanTransformer3DModel
[[autodoc]] WanTransformer3DModel
## Transformer2DModelOutput
[[autodoc]] models.modeling_outputs.Transformer2DModelOutput

View File

@@ -29,43 +29,3 @@ Customized normalization layers for supporting various models in 🤗 Diffusers.
## AdaGroupNorm
[[autodoc]] models.normalization.AdaGroupNorm
## AdaLayerNormContinuous
[[autodoc]] models.normalization.AdaLayerNormContinuous
## RMSNorm
[[autodoc]] models.normalization.RMSNorm
## GlobalResponseNorm
[[autodoc]] models.normalization.GlobalResponseNorm
## LuminaLayerNormContinuous
[[autodoc]] models.normalization.LuminaLayerNormContinuous
## SD35AdaLayerNormZeroX
[[autodoc]] models.normalization.SD35AdaLayerNormZeroX
## AdaLayerNormZeroSingle
[[autodoc]] models.normalization.AdaLayerNormZeroSingle
## LuminaRMSNormZero
[[autodoc]] models.normalization.LuminaRMSNormZero
## LpNorm
[[autodoc]] models.normalization.LpNorm
## CogView3PlusAdaLayerNormZeroTextImage
[[autodoc]] models.normalization.CogView3PlusAdaLayerNormZeroTextImage
## CogVideoXLayerNormZero
[[autodoc]] models.normalization.CogVideoXLayerNormZero
## MochiRMSNormZero
[[autodoc]] models.transformers.transformer_mochi.MochiRMSNormZero
## MochiRMSNorm
[[autodoc]] models.normalization.MochiRMSNorm

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text-to-Video Generation with AnimateDiff
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
## Overview
[AnimateDiff: Animate Your Personalized Text-to-Image Diffusion Models without Specific Tuning](https://arxiv.org/abs/2307.04725) by Yuwei Guo, Ceyuan Yang, Anyi Rao, Yaohui Wang, Yu Qiao, Dahua Lin, Bo Dai.
@@ -966,7 +962,7 @@ pipe.to("cuda")
prompt = {
0: "A caterpillar on a leaf, high quality, photorealistic",
40: "A caterpillar transforming into a cocoon, on a leaf, near flowers, photorealistic",
80: "A cocoon on a leaf, flowers in the background, photorealistic",
80: "A cocoon on a leaf, flowers in the backgrond, photorealistic",
120: "A cocoon maturing and a butterfly being born, flowers and leaves visible in the background, photorealistic",
160: "A beautiful butterfly, vibrant colors, sitting on a leaf, flowers in the background, photorealistic",
200: "A beautiful butterfly, flying away in a forest, photorealistic",

View File

@@ -62,50 +62,6 @@ image = pipeline(prompt).images[0]
image.save("auraflow.png")
```
Loading [GGUF checkpoints](https://huggingface.co/docs/diffusers/quantization/gguf) are also supported:
```py
import torch
from diffusers import (
AuraFlowPipeline,
GGUFQuantizationConfig,
AuraFlowTransformer2DModel,
)
transformer = AuraFlowTransformer2DModel.from_single_file(
"https://huggingface.co/city96/AuraFlow-v0.3-gguf/blob/main/aura_flow_0.3-Q2_K.gguf",
quantization_config=GGUFQuantizationConfig(compute_dtype=torch.bfloat16),
torch_dtype=torch.bfloat16,
)
pipeline = AuraFlowPipeline.from_pretrained(
"fal/AuraFlow-v0.3",
transformer=transformer,
torch_dtype=torch.bfloat16,
)
prompt = "a cute pony in a field of flowers"
image = pipeline(prompt).images[0]
image.save("auraflow.png")
```
## Support for `torch.compile()`
AuraFlow can be compiled with `torch.compile()` to speed up inference latency even for different resolutions. First, install PyTorch nightly following the instructions from [here](https://pytorch.org/). The snippet below shows the changes needed to enable this:
```diff
+ torch.fx.experimental._config.use_duck_shape = False
+ pipeline.transformer = torch.compile(
pipeline.transformer, fullgraph=True, dynamic=True
)
```
Specifying `use_duck_shape` to be `False` instructs the compiler if it should use the same symbolic variable to represent input sizes that are the same. For more details, check out [this comment](https://github.com/huggingface/diffusers/pull/11327#discussion_r2047659790).
This enables from 100% (on low resolutions) to a 30% (on 1536x1536 resolution) speed improvements.
Thanks to [AstraliteHeart](https://github.com/huggingface/diffusers/pull/11297/) who helped us rewrite the [`AuraFlowTransformer2DModel`] class so that the above works for different resolutions ([PR](https://github.com/huggingface/diffusers/pull/11297/)).
## AuraFlowPipeline
[[autodoc]] AuraFlowPipeline

View File

@@ -15,10 +15,6 @@
# CogVideoX
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer](https://arxiv.org/abs/2408.06072) from Tsinghua University & ZhipuAI, by Zhuoyi Yang, Jiayan Teng, Wendi Zheng, Ming Ding, Shiyu Huang, Jiazheng Xu, Yuanming Yang, Wenyi Hong, Xiaohan Zhang, Guanyu Feng, Da Yin, Xiaotao Gu, Yuxuan Zhang, Weihan Wang, Yean Cheng, Ting Liu, Bin Xu, Yuxiao Dong, Jie Tang.
The abstract from the paper is:

View File

@@ -1,34 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# CogView4
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [zRzRzRzRzRzRzR](https://github.com/zRzRzRzRzRzRzR). The original codebase can be found [here](https://huggingface.co/THUDM). The original weights can be found under [hf.co/THUDM](https://huggingface.co/THUDM).
## CogView4Pipeline
[[autodoc]] CogView4Pipeline
- all
- __call__
## CogView4PipelineOutput
[[autodoc]] pipelines.cogview4.pipeline_output.CogView4PipelineOutput

View File

@@ -1,64 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# ConsisID
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Identity-Preserving Text-to-Video Generation by Frequency Decomposition](https://arxiv.org/abs/2411.17440) from Peking University & University of Rochester & etc, by Shenghai Yuan, Jinfa Huang, Xianyi He, Yunyang Ge, Yujun Shi, Liuhan Chen, Jiebo Luo, Li Yuan.
The abstract from the paper is:
*Identity-preserving text-to-video (IPT2V) generation aims to create high-fidelity videos with consistent human identity. It is an important task in video generation but remains an open problem for generative models. This paper pushes the technical frontier of IPT2V in two directions that have not been resolved in the literature: (1) A tuning-free pipeline without tedious case-by-case finetuning, and (2) A frequency-aware heuristic identity-preserving Diffusion Transformer (DiT)-based control scheme. To achieve these goals, we propose **ConsisID**, a tuning-free DiT-based controllable IPT2V model to keep human-**id**entity **consis**tent in the generated video. Inspired by prior findings in frequency analysis of vision/diffusion transformers, it employs identity-control signals in the frequency domain, where facial features can be decomposed into low-frequency global features (e.g., profile, proportions) and high-frequency intrinsic features (e.g., identity markers that remain unaffected by pose changes). First, from a low-frequency perspective, we introduce a global facial extractor, which encodes the reference image and facial key points into a latent space, generating features enriched with low-frequency information. These features are then integrated into the shallow layers of the network to alleviate training challenges associated with DiT. Second, from a high-frequency perspective, we design a local facial extractor to capture high-frequency details and inject them into the transformer blocks, enhancing the model's ability to preserve fine-grained features. To leverage the frequency information for identity preservation, we propose a hierarchical training strategy, transforming a vanilla pre-trained video generation model into an IPT2V model. Extensive experiments demonstrate that our frequency-aware heuristic scheme provides an optimal control solution for DiT-based models. Thanks to this scheme, our **ConsisID** achieves excellent results in generating high-quality, identity-preserving videos, making strides towards more effective IPT2V. The model weight of ConsID is publicly available at https://github.com/PKU-YuanGroup/ConsisID.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [SHYuanBest](https://github.com/SHYuanBest). The original codebase can be found [here](https://github.com/PKU-YuanGroup/ConsisID). The original weights can be found under [hf.co/BestWishYsh](https://huggingface.co/BestWishYsh).
There are two official ConsisID checkpoints for identity-preserving text-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`BestWishYsh/ConsisID-preview`](https://huggingface.co/BestWishYsh/ConsisID-preview) | torch.bfloat16 |
| [`BestWishYsh/ConsisID-1.5`](https://huggingface.co/BestWishYsh/ConsisID-preview) | torch.bfloat16 |
### Memory optimization
ConsisID requires about 44 GB of GPU memory to decode 49 frames (6 seconds of video at 8 FPS) with output resolution 720x480 (W x H), which makes it not possible to run on consumer GPUs or free-tier T4 Colab. The following memory optimizations could be used to reduce the memory footprint. For replication, you can refer to [this](https://gist.github.com/SHYuanBest/bc4207c36f454f9e969adbb50eaf8258) script.
| Feature (overlay the previous) | Max Memory Allocated | Max Memory Reserved |
| :----------------------------- | :------------------- | :------------------ |
| - | 37 GB | 44 GB |
| enable_model_cpu_offload | 22 GB | 25 GB |
| enable_sequential_cpu_offload | 16 GB | 22 GB |
| vae.enable_slicing | 16 GB | 22 GB |
| vae.enable_tiling | 5 GB | 7 GB |
## ConsisIDPipeline
[[autodoc]] ConsisIDPipeline
- all
- __call__
## ConsisIDPipelineOutput
[[autodoc]] pipelines.consisid.pipeline_output.ConsisIDPipelineOutput

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# FluxControlInpaint
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
FluxControlInpaintPipeline is an implementation of Inpainting for Flux.1 Depth/Canny models. It is a pipeline that allows you to inpaint images using the Flux.1 Depth/Canny models. The pipeline takes an image and a mask as input and returns the inpainted image.
FLUX.1 Depth and Canny [dev] is a 12 billion parameter rectified flow transformer capable of generating an image based on a text description while following the structure of a given input image. **This is not a ControlNet model**.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet with Flux.1
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
FluxControlNetPipeline is an implementation of ControlNet for Flux.1.
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.

View File

@@ -1,36 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# ControlNet
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
The abstract from the paper is:
*We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.*
This pipeline was contributed by [ishan24](https://huggingface.co/ishan24). ❤️
The original codebase can be found at [NVlabs/Sana](https://github.com/NVlabs/Sana), and you can find official ControlNet checkpoints on [Efficient-Large-Model's](https://huggingface.co/Efficient-Large-Model) Hub profile.
## SanaControlNetPipeline
[[autodoc]] SanaControlNetPipeline
- all
- __call__
## SanaPipelineOutput
[[autodoc]] pipelines.sana.pipeline_output.SanaPipelineOutput

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet with Stable Diffusion 3
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
StableDiffusion3ControlNetPipeline is an implementation of ControlNet for Stable Diffusion 3.
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet with Stable Diffusion XL
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet was introduced in [Adding Conditional Control to Text-to-Image Diffusion Models](https://huggingface.co/papers/2302.05543) by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNetUnion
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNetUnionModel is an implementation of ControlNet for Stable Diffusion XL.
The ControlNet model was introduced in [ControlNetPlus](https://github.com/xinsir6/ControlNetPlus) by xinsir6. It supports multiple conditioning inputs without increasing computation.

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# ControlNet-XS
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
ControlNet-XS was introduced in [ControlNet-XS](https://vislearn.github.io/ControlNet-XS/) by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the [original ControlNet](https://huggingface.co/papers/2302.05543) can be made much smaller and still produce good results.
Like the original ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.

View File

@@ -1,41 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# Cosmos
[Cosmos World Foundation Model Platform for Physical AI](https://huggingface.co/papers/2501.03575) by NVIDIA.
*Physical AI needs to be trained digitally first. It needs a digital twin of itself, the policy model, and a digital twin of the world, the world model. In this paper, we present the Cosmos World Foundation Model Platform to help developers build customized world models for their Physical AI setups. We position a world foundation model as a general-purpose world model that can be fine-tuned into customized world models for downstream applications. Our platform covers a video curation pipeline, pre-trained world foundation models, examples of post-training of pre-trained world foundation models, and video tokenizers. To help Physical AI builders solve the most critical problems of our society, we make our platform open-source and our models open-weight with permissive licenses available via https://github.com/NVIDIA/Cosmos.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## CosmosTextToWorldPipeline
[[autodoc]] CosmosTextToWorldPipeline
- all
- __call__
## CosmosVideoToWorldPipeline
[[autodoc]] CosmosVideoToWorldPipeline
- all
- __call__
## CosmosPipelineOutput
[[autodoc]] pipelines.cosmos.pipeline_output.CosmosPipelineOutput

View File

@@ -12,11 +12,6 @@ specific language governing permissions and limitations under the License.
# DeepFloyd IF
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
## Overview
DeepFloyd IF is a novel state-of-the-art open-source text-to-image model with a high degree of photorealism and language understanding.

View File

@@ -1,88 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# EasyAnimate
[EasyAnimate](https://github.com/aigc-apps/EasyAnimate) by Alibaba PAI.
The description from it's GitHub page:
*EasyAnimate is a pipeline based on the transformer architecture, designed for generating AI images and videos, and for training baseline models and Lora models for Diffusion Transformer. We support direct prediction from pre-trained EasyAnimate models, allowing for the generation of videos with various resolutions, approximately 6 seconds in length, at 8fps (EasyAnimateV5.1, 1 to 49 frames). Additionally, users can train their own baseline and Lora models for specific style transformations.*
This pipeline was contributed by [bubbliiiing](https://github.com/bubbliiiing). The original codebase can be found [here](https://huggingface.co/alibaba-pai). The original weights can be found under [hf.co/alibaba-pai](https://huggingface.co/alibaba-pai).
There are two official EasyAnimate checkpoints for text-to-video and video-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`alibaba-pai/EasyAnimateV5.1-12b-zh`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh) | torch.float16 |
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-InP`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-InP) | torch.float16 |
There is one official EasyAnimate checkpoints available for image-to-video and video-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-InP`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-InP) | torch.float16 |
There are two official EasyAnimate checkpoints available for control-to-video.
| checkpoints | recommended inference dtype |
|:---:|:---:|
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-Control`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-Control) | torch.float16 |
| [`alibaba-pai/EasyAnimateV5.1-12b-zh-Control-Camera`](https://huggingface.co/alibaba-pai/EasyAnimateV5.1-12b-zh-Control-Camera) | torch.float16 |
For the EasyAnimateV5.1 series:
- Text-to-video (T2V) and Image-to-video (I2V) works for multiple resolutions. The width and height can vary from 256 to 1024.
- Both T2V and I2V models support generation with 1~49 frames and work best at this value. Exporting videos at 8 FPS is recommended.
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`EasyAnimatePipeline`] for inference with bitsandbytes.
```py
import torch
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, EasyAnimateTransformer3DModel, EasyAnimatePipeline
from diffusers.utils import export_to_video
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True)
transformer_8bit = EasyAnimateTransformer3DModel.from_pretrained(
"alibaba-pai/EasyAnimateV5.1-12b-zh",
subfolder="transformer",
quantization_config=quant_config,
torch_dtype=torch.float16,
)
pipeline = EasyAnimatePipeline.from_pretrained(
"alibaba-pai/EasyAnimateV5.1-12b-zh",
transformer=transformer_8bit,
torch_dtype=torch.float16,
device_map="balanced",
)
prompt = "A cat walks on the grass, realistic style."
negative_prompt = "bad detailed"
video = pipeline(prompt=prompt, negative_prompt=negative_prompt, num_frames=49, num_inference_steps=30).frames[0]
export_to_video(video, "cat.mp4", fps=8)
```
## EasyAnimatePipeline
[[autodoc]] EasyAnimatePipeline
- all
- __call__
## EasyAnimatePipelineOutput
[[autodoc]] pipelines.easyanimate.pipeline_output.EasyAnimatePipelineOutput

View File

@@ -12,11 +12,6 @@ specific language governing permissions and limitations under the License.
# Flux
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
Flux is a series of text-to-image generation models based on diffusion transformers. To know more about Flux, check out the original [blog post](https://blackforestlabs.ai/announcing-black-forest-labs/) by the creators of Flux, Black Forest Labs.
Original model checkpoints for Flux can be found [here](https://huggingface.co/black-forest-labs). Original inference code can be found [here](https://github.com/black-forest-labs/flux).
@@ -314,120 +309,7 @@ image.save("output.png")
When unloading the Control LoRA weights, call `pipe.unload_lora_weights(reset_to_overwritten_params=True)` to reset the `pipe.transformer` completely back to its original form. The resultant pipeline can then be used with methods like [`DiffusionPipeline.from_pipe`]. More details about this argument are available in [this PR](https://github.com/huggingface/diffusers/pull/10397).
## IP-Adapter
<Tip>
Check out [IP-Adapter](../../../using-diffusers/ip_adapter) to learn more about how IP-Adapters work.
</Tip>
An IP-Adapter lets you prompt Flux with images, in addition to the text prompt. This is especially useful when describing complex concepts that are difficult to articulate through text alone and you have reference images.
```python
import torch
from diffusers import FluxPipeline
from diffusers.utils import load_image
pipe = FluxPipeline.from_pretrained(
"black-forest-labs/FLUX.1-dev", torch_dtype=torch.bfloat16
).to("cuda")
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flux_ip_adapter_input.jpg").resize((1024, 1024))
pipe.load_ip_adapter(
"XLabs-AI/flux-ip-adapter",
weight_name="ip_adapter.safetensors",
image_encoder_pretrained_model_name_or_path="openai/clip-vit-large-patch14"
)
pipe.set_ip_adapter_scale(1.0)
image = pipe(
width=1024,
height=1024,
prompt="wearing sunglasses",
negative_prompt="",
true_cfg_scale=4.0,
generator=torch.Generator().manual_seed(4444),
ip_adapter_image=image,
).images[0]
image.save('flux_ip_adapter_output.jpg')
```
<div class="justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flux_ip_adapter_output.jpg"/>
<figcaption class="mt-2 text-sm text-center text-gray-500">IP-Adapter examples with prompt "wearing sunglasses"</figcaption>
</div>
## Optimize
Flux is a very large model and requires ~50GB of RAM/VRAM to load all the modeling components. Enable some of the optimizations below to lower the memory requirements.
### Group offloading
[Group offloading](../../optimization/memory#group-offloading) lowers VRAM usage by offloading groups of internal layers rather than the whole model or weights. You need to use [`~hooks.apply_group_offloading`] on all the model components of a pipeline. The `offload_type` parameter allows you to toggle between block and leaf-level offloading. Setting it to `leaf_level` offloads the lowest leaf-level parameters to the CPU instead of offloading at the module-level.
On CUDA devices that support asynchronous data streaming, set `use_stream=True` to overlap data transfer and computation to accelerate inference.
> [!TIP]
> It is possible to mix block and leaf-level offloading for different components in a pipeline.
```py
import torch
from diffusers import FluxPipeline
from diffusers.hooks import apply_group_offloading
model_id = "black-forest-labs/FLUX.1-dev"
dtype = torch.bfloat16
pipe = FluxPipeline.from_pretrained(
model_id,
torch_dtype=dtype,
)
apply_group_offloading(
pipe.transformer,
offload_type="leaf_level",
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
use_stream=True,
)
apply_group_offloading(
pipe.text_encoder,
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
offload_type="leaf_level",
use_stream=True,
)
apply_group_offloading(
pipe.text_encoder_2,
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
offload_type="leaf_level",
use_stream=True,
)
apply_group_offloading(
pipe.vae,
offload_device=torch.device("cpu"),
onload_device=torch.device("cuda"),
offload_type="leaf_level",
use_stream=True,
)
prompt="A cat wearing sunglasses and working as a lifeguard at pool."
generator = torch.Generator().manual_seed(181201)
image = pipe(
prompt,
width=576,
height=1024,
num_inference_steps=30,
generator=generator
).images[0]
image
```
### Running FP16 inference
## Running FP16 inference
Flux can generate high-quality images with FP16 (i.e. to accelerate inference on Turing/Volta GPUs) but produces different outputs compared to FP32/BF16. The issue is that some activations in the text encoders have to be clipped when running in FP16, which affects the overall image. Forcing text encoders to run with FP32 inference thus removes this output difference. See [here](https://github.com/huggingface/diffusers/pull/9097#issuecomment-2272292516) for details.
@@ -456,7 +338,7 @@ out = pipe(
out.save("image.png")
```
### Quantization
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
@@ -485,7 +367,7 @@ transformer_8bit = FluxTransformer2DModel.from_pretrained(
pipeline = FluxPipeline.from_pretrained(
"black-forest-labs/FLUX.1-dev",
text_encoder_2=text_encoder_8bit,
text_encoder=text_encoder_8bit,
transformer=transformer_8bit,
torch_dtype=torch.float16,
device_map="balanced",

View File

@@ -1,209 +0,0 @@
<!-- Copyright 2025 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# Framepack
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Packing Input Frame Context in Next-Frame Prediction Models for Video Generation](https://arxiv.org/abs/2504.12626) by Lvmin Zhang and Maneesh Agrawala.
*We present a neural network structure, FramePack, to train next-frame (or next-frame-section) prediction models for video generation. The FramePack compresses input frames to make the transformer context length a fixed number regardless of the video length. As a result, we are able to process a large number of frames using video diffusion with computation bottleneck similar to image diffusion. This also makes the training video batch sizes significantly higher (batch sizes become comparable to image diffusion training). We also propose an anti-drifting sampling method that generates frames in inverted temporal order with early-established endpoints to avoid exposure bias (error accumulation over iterations). Finally, we show that existing video diffusion models can be finetuned with FramePack, and their visual quality may be improved because the next-frame prediction supports more balanced diffusion schedulers with less extreme flow shift timesteps.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Available models
| Model name | Description |
|:---|:---|
- [`lllyasviel/FramePackI2V_HY`](https://huggingface.co/lllyasviel/FramePackI2V_HY) | Trained with the "inverted anti-drifting" strategy as described in the paper. Inference requires setting `sampling_type="inverted_anti_drifting"` when running the pipeline. |
- [`lllyasviel/FramePack_F1_I2V_HY_20250503`](https://huggingface.co/lllyasviel/FramePack_F1_I2V_HY_20250503) | Trained with a novel anti-drifting strategy but inference is performed in "vanilla" strategy as described in the paper. Inference requires setting `sampling_type="vanilla"` when running the pipeline. |
## Usage
Refer to the pipeline documentation for basic usage examples. The following section contains examples of offloading, different sampling methods, quantization, and more.
### First and last frame to video
The following example shows how to use Framepack with start and end image controls, using the inverted anti-drifiting sampling model.
```python
import torch
from diffusers import HunyuanVideoFramepackPipeline, HunyuanVideoFramepackTransformer3DModel
from diffusers.utils import export_to_video, load_image
from transformers import SiglipImageProcessor, SiglipVisionModel
transformer = HunyuanVideoFramepackTransformer3DModel.from_pretrained(
"lllyasviel/FramePackI2V_HY", torch_dtype=torch.bfloat16
)
feature_extractor = SiglipImageProcessor.from_pretrained(
"lllyasviel/flux_redux_bfl", subfolder="feature_extractor"
)
image_encoder = SiglipVisionModel.from_pretrained(
"lllyasviel/flux_redux_bfl", subfolder="image_encoder", torch_dtype=torch.float16
)
pipe = HunyuanVideoFramepackPipeline.from_pretrained(
"hunyuanvideo-community/HunyuanVideo",
transformer=transformer,
feature_extractor=feature_extractor,
image_encoder=image_encoder,
torch_dtype=torch.float16,
)
# Enable memory optimizations
pipe.enable_model_cpu_offload()
pipe.vae.enable_tiling()
prompt = "CG animation style, a small blue bird takes off from the ground, flapping its wings. The bird's feathers are delicate, with a unique pattern on its chest. The background shows a blue sky with white clouds under bright sunshine. The camera follows the bird upward, capturing its flight and the vastness of the sky from a close-up, low-angle perspective."
first_image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flf2v_input_first_frame.png"
)
last_image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flf2v_input_last_frame.png"
)
output = pipe(
image=first_image,
last_image=last_image,
prompt=prompt,
height=512,
width=512,
num_frames=91,
num_inference_steps=30,
guidance_scale=9.0,
generator=torch.Generator().manual_seed(0),
sampling_type="inverted_anti_drifting",
).frames[0]
export_to_video(output, "output.mp4", fps=30)
```
### Vanilla sampling
The following example shows how to use Framepack with the F1 model trained with vanilla sampling but new regulation approach for anti-drifting.
```python
import torch
from diffusers import HunyuanVideoFramepackPipeline, HunyuanVideoFramepackTransformer3DModel
from diffusers.utils import export_to_video, load_image
from transformers import SiglipImageProcessor, SiglipVisionModel
transformer = HunyuanVideoFramepackTransformer3DModel.from_pretrained(
"lllyasviel/FramePack_F1_I2V_HY_20250503", torch_dtype=torch.bfloat16
)
feature_extractor = SiglipImageProcessor.from_pretrained(
"lllyasviel/flux_redux_bfl", subfolder="feature_extractor"
)
image_encoder = SiglipVisionModel.from_pretrained(
"lllyasviel/flux_redux_bfl", subfolder="image_encoder", torch_dtype=torch.float16
)
pipe = HunyuanVideoFramepackPipeline.from_pretrained(
"hunyuanvideo-community/HunyuanVideo",
transformer=transformer,
feature_extractor=feature_extractor,
image_encoder=image_encoder,
torch_dtype=torch.float16,
)
# Enable memory optimizations
pipe.enable_model_cpu_offload()
pipe.vae.enable_tiling()
image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/penguin.png"
)
output = pipe(
image=image,
prompt="A penguin dancing in the snow",
height=832,
width=480,
num_frames=91,
num_inference_steps=30,
guidance_scale=9.0,
generator=torch.Generator().manual_seed(0),
sampling_type="vanilla",
).frames[0]
export_to_video(output, "output.mp4", fps=30)
```
### Group offloading
Group offloading ([`~hooks.apply_group_offloading`]) provides aggressive memory optimizations for offloading internal parts of any model to the CPU, with possibly no additional overhead to generation time. If you have very low VRAM available, this approach may be suitable for you depending on the amount of CPU RAM available.
```python
import torch
from diffusers import HunyuanVideoFramepackPipeline, HunyuanVideoFramepackTransformer3DModel
from diffusers.hooks import apply_group_offloading
from diffusers.utils import export_to_video, load_image
from transformers import SiglipImageProcessor, SiglipVisionModel
transformer = HunyuanVideoFramepackTransformer3DModel.from_pretrained(
"lllyasviel/FramePack_F1_I2V_HY_20250503", torch_dtype=torch.bfloat16
)
feature_extractor = SiglipImageProcessor.from_pretrained(
"lllyasviel/flux_redux_bfl", subfolder="feature_extractor"
)
image_encoder = SiglipVisionModel.from_pretrained(
"lllyasviel/flux_redux_bfl", subfolder="image_encoder", torch_dtype=torch.float16
)
pipe = HunyuanVideoFramepackPipeline.from_pretrained(
"hunyuanvideo-community/HunyuanVideo",
transformer=transformer,
feature_extractor=feature_extractor,
image_encoder=image_encoder,
torch_dtype=torch.float16,
)
# Enable group offloading
onload_device = torch.device("cuda")
offload_device = torch.device("cpu")
list(map(
lambda x: apply_group_offloading(x, onload_device, offload_device, offload_type="leaf_level", use_stream=True, low_cpu_mem_usage=True),
[pipe.text_encoder, pipe.text_encoder_2, pipe.transformer]
))
pipe.image_encoder.to(onload_device)
pipe.vae.to(onload_device)
pipe.vae.enable_tiling()
image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/penguin.png"
)
output = pipe(
image=image,
prompt="A penguin dancing in the snow",
height=832,
width=480,
num_frames=91,
num_inference_steps=30,
guidance_scale=9.0,
generator=torch.Generator().manual_seed(0),
sampling_type="vanilla",
).frames[0]
print(f"Max memory: {torch.cuda.max_memory_allocated() / 1024**3:.3f} GB")
export_to_video(output, "output.mp4", fps=30)
```
## HunyuanVideoFramepackPipeline
[[autodoc]] HunyuanVideoFramepackPipeline
- all
- __call__
## HunyuanVideoPipelineOutput
[[autodoc]] pipelines.hunyuan_video.pipeline_output.HunyuanVideoPipelineOutput

View File

@@ -1,43 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# HiDreamImage
[HiDream-I1](https://huggingface.co/HiDream-ai) by HiDream.ai
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Available models
The following models are available for the [`HiDreamImagePipeline`](text-to-image) pipeline:
| Model name | Description |
|:---|:---|
| [`HiDream-ai/HiDream-I1-Full`](https://huggingface.co/HiDream-ai/HiDream-I1-Full) | - |
| [`HiDream-ai/HiDream-I1-Dev`](https://huggingface.co/HiDream-ai/HiDream-I1-Dev) | - |
| [`HiDream-ai/HiDream-I1-Fast`](https://huggingface.co/HiDream-ai/HiDream-I1-Fast) | - |
## HiDreamImagePipeline
[[autodoc]] HiDreamImagePipeline
- all
- __call__
## HiDreamImagePipelineOutput
[[autodoc]] pipelines.hidream_image.pipeline_output.HiDreamImagePipelineOutput

View File

@@ -14,13 +14,9 @@
# HunyuanVideo
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[HunyuanVideo](https://www.arxiv.org/abs/2412.03603) by Tencent.
*Recent advancements in video generation have significantly impacted daily life for both individuals and industries. However, the leading video generation models remain closed-source, resulting in a notable performance gap between industry capabilities and those available to the public. In this report, we introduce HunyuanVideo, an innovative open-source video foundation model that demonstrates performance in video generation comparable to, or even surpassing, that of leading closed-source models. HunyuanVideo encompasses a comprehensive framework that integrates several key elements, including data curation, advanced architectural design, progressive model scaling and training, and an efficient infrastructure tailored for large-scale model training and inference. As a result, we successfully trained a video generative model with over 13 billion parameters, making it the largest among all open-source models. We conducted extensive experiments and implemented a series of targeted designs to ensure high visual quality, motion dynamics, text-video alignment, and advanced filming techniques. According to evaluations by professionals, HunyuanVideo outperforms previous state-of-the-art models, including Runway Gen-3, Luma 1.6, and three top-performing Chinese video generative models. By releasing the code for the foundation model and its applications, we aim to bridge the gap between closed-source and open-source communities. This initiative will empower individuals within the community to experiment with their ideas, fostering a more dynamic and vibrant video generation ecosystem. The code is publicly available at [this https URL](https://github.com/tencent/HunyuanVideo).*
*Recent advancements in video generation have significantly impacted daily life for both individuals and industries. However, the leading video generation models remain closed-source, resulting in a notable performance gap between industry capabilities and those available to the public. In this report, we introduce HunyuanVideo, an innovative open-source video foundation model that demonstrates performance in video generation comparable to, or even surpassing, that of leading closed-source models. HunyuanVideo encompasses a comprehensive framework that integrates several key elements, including data curation, advanced architectural design, progressive model scaling and training, and an efficient infrastructure tailored for large-scale model training and inference. As a result, we successfully trained a video generative model with over 13 billion parameters, making it the largest among all open-source models. We conducted extensive experiments and implemented a series of targeted designs to ensure high visual quality, motion dynamics, text-video alignment, and advanced filming techniques. According to evaluations by professionals, HunyuanVideo outperforms previous state-of-the-art models, including Runway Gen-3, Luma 1.6, and three top-performing Chinese video generative models. By releasing the code for the foundation model and its applications, we aim to bridge the gap between closed-source and open-source communities. This initiative will empower individuals within the community to experiment with their ideas, fostering a more dynamic and vibrant video generation ecosystem. The code is publicly available at [this https URL](https://github.com/Tencent/HunyuanVideo).*
<Tip>
@@ -36,23 +32,6 @@ Recommendations for inference:
- For smaller resolution videos, try lower values of `shift` (between `2.0` to `5.0`) in the [Scheduler](https://huggingface.co/docs/diffusers/main/en/api/schedulers/flow_match_euler_discrete#diffusers.FlowMatchEulerDiscreteScheduler.shift). For larger resolution images, try higher values (between `7.0` and `12.0`). The default value is `7.0` for HunyuanVideo.
- For more information about supported resolutions and other details, please refer to the original repository [here](https://github.com/Tencent/HunyuanVideo/).
## Available models
The following models are available for the [`HunyuanVideoPipeline`](text-to-video) pipeline:
| Model name | Description |
|:---|:---|
| [`hunyuanvideo-community/HunyuanVideo`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo) | Official HunyuanVideo (guidance-distilled). Performs best at multiple resolutions and frames. Performs best with `guidance_scale=6.0`, `true_cfg_scale=1.0` and without a negative prompt. |
| [`https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-T2V`](https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-T2V) | Skywork's custom finetune of HunyuanVideo (de-distilled). Performs best with `97x544x960` resolution, `guidance_scale=1.0`, `true_cfg_scale=6.0` and a negative prompt. |
The following models are available for the image-to-video pipeline:
| Model name | Description |
|:---|:---|
| [`Skywork/SkyReels-V1-Hunyuan-I2V`](https://huggingface.co/Skywork/SkyReels-V1-Hunyuan-I2V) | Skywork's custom finetune of HunyuanVideo (de-distilled). Performs best with `97x544x960` resolution. Performs best at `97x544x960` resolution, `guidance_scale=1.0`, `true_cfg_scale=6.0` and a negative prompt. |
| [`hunyuanvideo-community/HunyuanVideo-I2V-33ch`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo-I2V) | Tecent's official HunyuanVideo 33-channel I2V model. Performs best at resolutions of 480, 720, 960, 1280. A higher `shift` value when initializing the scheduler is recommended (good values are between 7 and 20). |
| [`hunyuanvideo-community/HunyuanVideo-I2V`](https://huggingface.co/hunyuanvideo-community/HunyuanVideo-I2V) | Tecent's official HunyuanVideo 16-channel I2V model. Performs best at resolutions of 480, 720, 960, 1280. A higher `shift` value when initializing the scheduler is recommended (good values are between 7 and 20) |
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
@@ -66,14 +45,14 @@ from diffusers.utils import export_to_video
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True)
transformer_8bit = HunyuanVideoTransformer3DModel.from_pretrained(
"hunyuanvideo-community/HunyuanVideo",
"tencent/HunyuanVideo",
subfolder="transformer",
quantization_config=quant_config,
torch_dtype=torch.bfloat16,
torch_dtype=torch.float16,
)
pipeline = HunyuanVideoPipeline.from_pretrained(
"hunyuanvideo-community/HunyuanVideo",
"tencent/HunyuanVideo",
transformer=transformer_8bit,
torch_dtype=torch.float16,
device_map="balanced",

View File

@@ -9,10 +9,6 @@ specific language governing permissions and limitations under the License.
# Kandinsky 3
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Kandinsky 3 is created by [Vladimir Arkhipkin](https://github.com/oriBetelgeuse),[Anastasia Maltseva](https://github.com/NastyaMittseva),[Igor Pavlov](https://github.com/boomb0om),[Andrei Filatov](https://github.com/anvilarth),[Arseniy Shakhmatov](https://github.com/cene555),[Andrey Kuznetsov](https://github.com/kuznetsoffandrey),[Denis Dimitrov](https://github.com/denndimitrov), [Zein Shaheen](https://github.com/zeinsh)
The description from it's GitHub page:

View File

@@ -12,11 +12,6 @@ specific language governing permissions and limitations under the License.
# Kolors: Effective Training of Diffusion Model for Photorealistic Text-to-Image Synthesis
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
![](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/kolors/kolors_header_collage.png)
Kolors is a large-scale text-to-image generation model based on latent diffusion, developed by [the Kuaishou Kolors team](https://github.com/Kwai-Kolors/Kolors). Trained on billions of text-image pairs, Kolors exhibits significant advantages over both open-source and closed-source models in visual quality, complex semantic accuracy, and text rendering for both Chinese and English characters. Furthermore, Kolors supports both Chinese and English inputs, demonstrating strong performance in understanding and generating Chinese-specific content. For more details, please refer to this [technical report](https://github.com/Kwai-Kolors/Kolors/blob/master/imgs/Kolors_paper.pdf).

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Latent Consistency Models
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Latent Consistency Models (LCMs) were proposed in [Latent Consistency Models: Synthesizing High-Resolution Images with Few-Step Inference](https://huggingface.co/papers/2310.04378) by Simian Luo, Yiqin Tan, Longbo Huang, Jian Li, and Hang Zhao.
The abstract of the paper is as follows:

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# LEDITS++
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
LEDITS++ was proposed in [LEDITS++: Limitless Image Editing using Text-to-Image Models](https://huggingface.co/papers/2311.16711) by Manuel Brack, Felix Friedrich, Katharina Kornmeier, Linoy Tsaban, Patrick Schramowski, Kristian Kersting, Apolinário Passos.
The abstract from the paper is:
@@ -29,7 +25,7 @@ You can find additional information about LEDITS++ on the [project page](https:/
</Tip>
<Tip warning={true}>
Due to some backward compatibility issues with the current diffusers implementation of [`~schedulers.DPMSolverMultistepScheduler`] this implementation of LEdits++ can no longer guarantee perfect inversion.
Due to some backward compatability issues with the current diffusers implementation of [`~schedulers.DPMSolverMultistepScheduler`] this implementation of LEdits++ can no longer guarantee perfect inversion.
This issue is unlikely to have any noticeable effects on applied use-cases. However, we provide an alternative implementation that guarantees perfect inversion in a dedicated [GitHub repo](https://github.com/ml-research/ledits_pp).
</Tip>

View File

@@ -14,11 +14,6 @@
# LTX Video
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
[LTX Video](https://huggingface.co/Lightricks/LTX-Video) is the first DiT-based video generation model capable of generating high-quality videos in real-time. It produces 24 FPS videos at a 768x512 resolution faster than they can be watched. Trained on a large-scale dataset of diverse videos, the model generates high-resolution videos with realistic and varied content. We provide a model for both text-to-video as well as image + text-to-video usecases.
<Tip>
@@ -31,103 +26,11 @@ Available models:
| Model name | Recommended dtype |
|:-------------:|:-----------------:|
| [`LTX Video 2B 0.9.0`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.safetensors) | `torch.bfloat16` |
| [`LTX Video 2B 0.9.1`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.1.safetensors) | `torch.bfloat16` |
| [`LTX Video 2B 0.9.5`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.5.safetensors) | `torch.bfloat16` |
| [`LTX Video 13B 0.9.7`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltxv-13b-0.9.7-dev.safetensors) | `torch.bfloat16` |
| [`LTX Video Spatial Upscaler 0.9.7`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltxv-spatial-upscaler-0.9.7.safetensors) | `torch.bfloat16` |
| [`LTX Video 0.9.0`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.safetensors) | `torch.bfloat16` |
| [`LTX Video 0.9.1`](https://huggingface.co/Lightricks/LTX-Video/blob/main/ltx-video-2b-v0.9.1.safetensors) | `torch.bfloat16` |
Note: The recommended dtype is for the transformer component. The VAE and text encoders can be either `torch.float32`, `torch.bfloat16` or `torch.float16` but the recommended dtype is `torch.bfloat16` as used in the original repository.
## Recommended settings for generation
For the best results, it is recommended to follow the guidelines mentioned in the official LTX Video [repository](https://github.com/Lightricks/LTX-Video).
- Some variants of LTX Video are guidance-distilled. For guidance-distilled models, `guidance_scale` must be set to `1.0`. For any other models, `guidance_scale` should be set higher (e.g., `5.0`) for good generation quality.
- For variants with a timestep-aware VAE (LTXV 0.9.1 and above), it is recommended to set `decode_timestep` to `0.05` and `image_cond_noise_scale` to `0.025`.
- For variants that support interpolation between multiple conditioning images and videos (LTXV 0.9.5 and above), it is recommended to use similar looking images/videos for the best results. High divergence between the conditionings may lead to abrupt transitions in the generated video.
## Using LTX Video 13B 0.9.7
LTX Video 0.9.7 comes with a spatial latent upscaler and a 13B parameter transformer. The inference involves generating a low resolution video first, which is very fast, followed by upscaling and refining the generated video.
<!-- TODO(aryan): modify when official checkpoints are available -->
```python
import torch
from diffusers import LTXConditionPipeline, LTXLatentUpsamplePipeline
from diffusers.pipelines.ltx.pipeline_ltx_condition import LTXVideoCondition
from diffusers.utils import export_to_video, load_video
pipe = LTXConditionPipeline.from_pretrained("a-r-r-o-w/LTX-Video-0.9.7-diffusers", torch_dtype=torch.bfloat16)
pipe_upsample = LTXLatentUpsamplePipeline.from_pretrained("a-r-r-o-w/LTX-Video-0.9.7-Latent-Spatial-Upsampler-diffusers", vae=pipe.vae, torch_dtype=torch.bfloat16)
pipe.to("cuda")
pipe_upsample.to("cuda")
pipe.vae.enable_tiling()
def round_to_nearest_resolution_acceptable_by_vae(height, width):
height = height - (height % pipe.vae_temporal_compression_ratio)
width = width - (width % pipe.vae_temporal_compression_ratio)
return height, width
video = load_video(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/cosmos/cosmos-video2world-input-vid.mp4"
)[:21] # Use only the first 21 frames as conditioning
condition1 = LTXVideoCondition(video=video, frame_index=0)
prompt = "The video depicts a winding mountain road covered in snow, with a single vehicle traveling along it. The road is flanked by steep, rocky cliffs and sparse vegetation. The landscape is characterized by rugged terrain and a river visible in the distance. The scene captures the solitude and beauty of a winter drive through a mountainous region."
negative_prompt = "worst quality, inconsistent motion, blurry, jittery, distorted"
expected_height, expected_width = 768, 1152
downscale_factor = 2 / 3
num_frames = 161
# Part 1. Generate video at smaller resolution
# Text-only conditioning is also supported without the need to pass `conditions`
downscaled_height, downscaled_width = int(expected_height * downscale_factor), int(expected_width * downscale_factor)
downscaled_height, downscaled_width = round_to_nearest_resolution_acceptable_by_vae(downscaled_height, downscaled_width)
latents = pipe(
conditions=[condition1],
prompt=prompt,
negative_prompt=negative_prompt,
width=downscaled_width,
height=downscaled_height,
num_frames=num_frames,
num_inference_steps=30,
generator=torch.Generator().manual_seed(0),
output_type="latent",
).frames
# Part 2. Upscale generated video using latent upsampler with fewer inference steps
# The available latent upsampler upscales the height/width by 2x
upscaled_height, upscaled_width = downscaled_height * 2, downscaled_width * 2
upscaled_latents = pipe_upsample(
latents=latents,
output_type="latent"
).frames
# Part 3. Denoise the upscaled video with few steps to improve texture (optional, but recommended)
video = pipe(
conditions=[condition1],
prompt=prompt,
negative_prompt=negative_prompt,
width=upscaled_width,
height=upscaled_height,
num_frames=num_frames,
denoise_strength=0.4, # Effectively, 4 inference steps out of 10
num_inference_steps=10,
latents=upscaled_latents,
decode_timestep=0.05,
image_cond_noise_scale=0.025,
generator=torch.Generator().manual_seed(0),
output_type="pil",
).frames[0]
# Part 4. Downscale the video to the expected resolution
video = [frame.resize((expected_width, expected_height)) for frame in video]
export_to_video(video, "output.mp4", fps=24)
```
## Loading Single Files
Loading the original LTX Video checkpoints is also possible with [`~ModelMixin.from_single_file`]. We recommend using `from_single_file` for the Lightricks series of models, as they plan to release multiple models in the future in the single file format.
@@ -289,18 +192,6 @@ export_to_video(video, "ship.mp4", fps=24)
- all
- __call__
## LTXConditionPipeline
[[autodoc]] LTXConditionPipeline
- all
- __call__
## LTXLatentUpsamplePipeline
[[autodoc]] LTXLatentUpsamplePipeline
- all
- __call__
## LTXPipelineOutput
[[autodoc]] pipelines.ltx.pipeline_output.LTXPipelineOutput

View File

@@ -58,10 +58,10 @@ Use [`torch.compile`](https://huggingface.co/docs/diffusers/main/en/tutorials/fa
First, load the pipeline:
```python
from diffusers import LuminaPipeline
from diffusers import LuminaText2ImgPipeline
import torch
pipeline = LuminaPipeline.from_pretrained(
pipeline = LuminaText2ImgPipeline.from_pretrained(
"Alpha-VLLM/Lumina-Next-SFT-diffusers", torch_dtype=torch.bfloat16
).to("cuda")
```
@@ -86,11 +86,11 @@ image = pipeline(prompt="Upper body of a young woman in a Victorian-era outfit w
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`LuminaPipeline`] for inference with bitsandbytes.
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`LuminaText2ImgPipeline`] for inference with bitsandbytes.
```py
import torch
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, Transformer2DModel, LuminaPipeline
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, Transformer2DModel, LuminaText2ImgPipeline
from transformers import BitsAndBytesConfig as BitsAndBytesConfig, T5EncoderModel
quant_config = BitsAndBytesConfig(load_in_8bit=True)
@@ -109,7 +109,7 @@ transformer_8bit = Transformer2DModel.from_pretrained(
torch_dtype=torch.float16,
)
pipeline = LuminaPipeline.from_pretrained(
pipeline = LuminaText2ImgPipeline.from_pretrained(
"Alpha-VLLM/Lumina-Next-SFT-diffusers",
text_encoder=text_encoder_8bit,
transformer=transformer_8bit,
@@ -122,9 +122,9 @@ image = pipeline(prompt).images[0]
image.save("lumina.png")
```
## LuminaPipeline
## LuminaText2ImgPipeline
[[autodoc]] LuminaPipeline
[[autodoc]] LuminaText2ImgPipeline
- all
- __call__

View File

@@ -1,87 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# Lumina2
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Lumina Image 2.0: A Unified and Efficient Image Generative Model](https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0) is a 2 billion parameter flow-based diffusion transformer capable of generating diverse images from text descriptions.
The abstract from the paper is:
*We introduce Lumina-Image 2.0, an advanced text-to-image model that surpasses previous state-of-the-art methods across multiple benchmarks, while also shedding light on its potential to evolve into a generalist vision intelligence model. Lumina-Image 2.0 exhibits three key properties: (1) Unification it adopts a unified architecture that treats text and image tokens as a joint sequence, enabling natural cross-modal interactions and facilitating task expansion. Besides, since high-quality captioners can provide semantically better-aligned text-image training pairs, we introduce a unified captioning system, UniCaptioner, which generates comprehensive and precise captions for the model. This not only accelerates model convergence but also enhances prompt adherence, variable-length prompt handling, and task generalization via prompt templates. (2) Efficiency to improve the efficiency of the unified architecture, we develop a set of optimization techniques that improve semantic learning and fine-grained texture generation during training while incorporating inference-time acceleration strategies without compromising image quality. (3) Transparency we open-source all training details, code, and models to ensure full reproducibility, aiming to bridge the gap between well-resourced closed-source research teams and independent developers.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
## Using Single File loading with Lumina Image 2.0
Single file loading for Lumina Image 2.0 is available for the `Lumina2Transformer2DModel`
```python
import torch
from diffusers import Lumina2Transformer2DModel, Lumina2Pipeline
ckpt_path = "https://huggingface.co/Alpha-VLLM/Lumina-Image-2.0/blob/main/consolidated.00-of-01.pth"
transformer = Lumina2Transformer2DModel.from_single_file(
ckpt_path, torch_dtype=torch.bfloat16
)
pipe = Lumina2Pipeline.from_pretrained(
"Alpha-VLLM/Lumina-Image-2.0", transformer=transformer, torch_dtype=torch.bfloat16
)
pipe.enable_model_cpu_offload()
image = pipe(
"a cat holding a sign that says hello",
generator=torch.Generator("cpu").manual_seed(0),
).images[0]
image.save("lumina-single-file.png")
```
## Using GGUF Quantized Checkpoints with Lumina Image 2.0
GGUF Quantized checkpoints for the `Lumina2Transformer2DModel` can be loaded via `from_single_file` with the `GGUFQuantizationConfig`
```python
from diffusers import Lumina2Transformer2DModel, Lumina2Pipeline, GGUFQuantizationConfig
ckpt_path = "https://huggingface.co/calcuis/lumina-gguf/blob/main/lumina2-q4_0.gguf"
transformer = Lumina2Transformer2DModel.from_single_file(
ckpt_path,
quantization_config=GGUFQuantizationConfig(compute_dtype=torch.bfloat16),
torch_dtype=torch.bfloat16,
)
pipe = Lumina2Pipeline.from_pretrained(
"Alpha-VLLM/Lumina-Image-2.0", transformer=transformer, torch_dtype=torch.bfloat16
)
pipe.enable_model_cpu_offload()
image = pipe(
"a cat holding a sign that says hello",
generator=torch.Generator("cpu").manual_seed(0),
).images[0]
image.save("lumina-gguf.png")
```
## Lumina2Pipeline
[[autodoc]] Lumina2Pipeline
- all
- __call__

View File

@@ -1,6 +1,4 @@
<!--
Copyright 2023-2025 Marigold Team, ETH Zürich. All rights reserved.
Copyright 2024-2025 The HuggingFace Team. All rights reserved.
<!--Copyright 2024 Marigold authors and The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
@@ -12,120 +10,67 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Marigold Computer Vision
# Marigold Pipelines for Computer Vision Tasks
![marigold](https://marigoldmonodepth.github.io/images/teaser_collage_compressed.jpg)
Marigold was proposed in
[Repurposing Diffusion-Based Image Generators for Monocular Depth Estimation](https://huggingface.co/papers/2312.02145),
a CVPR 2024 Oral paper by
[Bingxin Ke](http://www.kebingxin.com/),
[Anton Obukhov](https://www.obukhov.ai/),
[Shengyu Huang](https://shengyuh.github.io/),
[Nando Metzger](https://nandometzger.github.io/),
[Rodrigo Caye Daudt](https://rcdaudt.github.io/), and
[Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
The core idea is to **repurpose the generative prior of Text-to-Image Latent Diffusion Models (LDMs) for traditional
computer vision tasks**.
This approach was explored by fine-tuning Stable Diffusion for **Monocular Depth Estimation**, as demonstrated in the
teaser above.
Marigold was proposed in [Repurposing Diffusion-Based Image Generators for Monocular Depth Estimation](https://huggingface.co/papers/2312.02145), a CVPR 2024 Oral paper by [Bingxin Ke](http://www.kebingxin.com/), [Anton Obukhov](https://www.obukhov.ai/), [Shengyu Huang](https://shengyuh.github.io/), [Nando Metzger](https://nandometzger.github.io/), [Rodrigo Caye Daudt](https://rcdaudt.github.io/), and [Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
The idea is to repurpose the rich generative prior of Text-to-Image Latent Diffusion Models (LDMs) for traditional computer vision tasks.
Initially, this idea was explored to fine-tune Stable Diffusion for Monocular Depth Estimation, as shown in the teaser above.
Later,
- [Tianfu Wang](https://tianfwang.github.io/) trained the first Latent Consistency Model (LCM) of Marigold, which unlocked fast single-step inference;
- [Kevin Qu](https://www.linkedin.com/in/kevin-qu-b3417621b/?locale=en_US) extended the approach to Surface Normals Estimation;
- [Anton Obukhov](https://www.obukhov.ai/) contributed the pipelines and documentation into diffusers (enabled and supported by [YiYi Xu](https://yiyixuxu.github.io/) and [Sayak Paul](https://sayak.dev/)).
Marigold was later extended in the follow-up paper,
[Marigold: Affordable Adaptation of Diffusion-Based Image Generators for Image Analysis](https://huggingface.co/papers/2312.02145),
authored by
[Bingxin Ke](http://www.kebingxin.com/),
[Kevin Qu](https://www.linkedin.com/in/kevin-qu-b3417621b/?locale=en_US),
[Tianfu Wang](https://tianfwang.github.io/),
[Nando Metzger](https://nandometzger.github.io/),
[Shengyu Huang](https://shengyuh.github.io/),
[Bo Li](https://www.linkedin.com/in/bobboli0202/),
[Anton Obukhov](https://www.obukhov.ai/), and
[Konrad Schindler](https://scholar.google.com/citations?user=FZuNgqIAAAAJ&hl=en).
This work expanded Marigold to support new modalities such as **Surface Normals** and **Intrinsic Image Decomposition**
(IID), introduced a training protocol for **Latent Consistency Models** (LCM), and demonstrated **High-Resolution** (HR)
processing capability.
The abstract from the paper is:
<Tip>
The early Marigold models (`v1-0` and earlier) were optimized for best results with at least 10 inference steps.
LCM models were later developed to enable high-quality inference in just 1 to 4 steps.
Marigold models `v1-1` and later use the DDIM scheduler to achieve optimal
results in as few as 1 to 4 steps.
</Tip>
*Monocular depth estimation is a fundamental computer vision task. Recovering 3D depth from a single image is geometrically ill-posed and requires scene understanding, so it is not surprising that the rise of deep learning has led to a breakthrough. The impressive progress of monocular depth estimators has mirrored the growth in model capacity, from relatively modest CNNs to large Transformer architectures. Still, monocular depth estimators tend to struggle when presented with images with unfamiliar content and layout, since their knowledge of the visual world is restricted by the data seen during training, and challenged by zero-shot generalization to new domains. This motivates us to explore whether the extensive priors captured in recent generative diffusion models can enable better, more generalizable depth estimation. We introduce Marigold, a method for affine-invariant monocular depth estimation that is derived from Stable Diffusion and retains its rich prior knowledge. The estimator can be fine-tuned in a couple of days on a single GPU using only synthetic training data. It delivers state-of-the-art performance across a wide range of datasets, including over 20% performance gains in specific cases. Project page: https://marigoldmonodepth.github.io.*
## Available Pipelines
Each pipeline is tailored for a specific computer vision task, processing an input RGB image and generating a
corresponding prediction.
Currently, the following computer vision tasks are implemented:
Each pipeline supports one Computer Vision task, which takes an input RGB image as input and produces a *prediction* of the modality of interest, such as a depth map of the input image.
Currently, the following tasks are implemented:
| Pipeline | Predicted Modalities | Demos |
|---------------------------------------------------------------------------------------------------------------------------------------------|------------------------------------------------------------------------------------------------------------------|:--------------------------------------------------------------------------------------------------------------------------------------------------:|
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-lcm), [Slow Original Demo (DDIM)](https://huggingface.co/spaces/prs-eth/marigold) |
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) | [Fast Demo (LCM)](https://huggingface.co/spaces/prs-eth/marigold-normals-lcm) |
| Pipeline | Recommended Model Checkpoints | Spaces (Interactive Apps) | Predicted Modalities |
|---------------------------------------------------------------------------------------------------------------------------------------------------|---------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|:------------------------------------------------------------------------------------:|------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
| [MarigoldDepthPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_depth.py) | [prs-eth/marigold-depth-v1-1](https://huggingface.co/prs-eth/marigold-depth-v1-1) | [Depth Estimation](https://huggingface.co/spaces/prs-eth/marigold) | [Depth](https://en.wikipedia.org/wiki/Depth_map), [Disparity](https://en.wikipedia.org/wiki/Binocular_disparity) |
| [MarigoldNormalsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_normals.py) | [prs-eth/marigold-normals-v1-1](https://huggingface.co/prs-eth/marigold-normals-v1-1) | [Surface Normals Estimation](https://huggingface.co/spaces/prs-eth/marigold-normals) | [Surface normals](https://en.wikipedia.org/wiki/Normal_mapping) |
| [MarigoldIntrinsicsPipeline](https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/marigold/pipeline_marigold_intrinsics.py) | [prs-eth/marigold-iid-appearance-v1-1](https://huggingface.co/prs-eth/marigold-iid-appearance-v1-1),<br>[prs-eth/marigold-iid-lighting-v1-1](https://huggingface.co/prs-eth/marigold-iid-lighting-v1-1) | [Intrinsic Image Decomposition](https://huggingface.co/spaces/prs-eth/marigold-iid) | [Albedo](https://en.wikipedia.org/wiki/Albedo), [Materials](https://www.n.aiq3d.com/wiki/roughnessmetalnessao-map), [Lighting](https://en.wikipedia.org/wiki/Diffuse_reflection) |
## Available Checkpoints
All original checkpoints are available under the [PRS-ETH](https://huggingface.co/prs-eth/) organization on Hugging Face.
They are designed for use with diffusers pipelines and the [original codebase](https://github.com/prs-eth/marigold), which can also be used to train
new model checkpoints.
The following is a summary of the recommended checkpoints, all of which produce reliable results with 1 to 4 steps.
| Checkpoint | Modality | Comment |
|-----------------------------------------------------------------------------------------------------|--------------|--------------------------------------------------------------------------------------------------------------------------------------------------------------------------------------|
| [prs-eth/marigold-depth-v1-1](https://huggingface.co/prs-eth/marigold-depth-v1-1) | Depth | Affine-invariant depth prediction assigns each pixel a value between 0 (near plane) and 1 (far plane), with both planes determined by the model during inference. |
| [prs-eth/marigold-normals-v0-1](https://huggingface.co/prs-eth/marigold-normals-v0-1) | Normals | The surface normals predictions are unit-length 3D vectors in the screen space camera, with values in the range from -1 to 1. |
| [prs-eth/marigold-iid-appearance-v1-1](https://huggingface.co/prs-eth/marigold-iid-appearance-v1-1) | Intrinsics | InteriorVerse decomposition is comprised of Albedo and two BRDF material properties: Roughness and Metallicity. |
| [prs-eth/marigold-iid-lighting-v1-1](https://huggingface.co/prs-eth/marigold-iid-lighting-v1-1) | Intrinsics | HyperSim decomposition of an image &nbsp\\(I\\)&nbsp is comprised of Albedo &nbsp\\(A\\), Diffuse shading &nbsp\\(S\\), and Non-diffuse residual &nbsp\\(R\\): &nbsp\\(I = A*S+R\\). |
The original checkpoints can be found under the [PRS-ETH](https://huggingface.co/prs-eth/) Hugging Face organization.
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff
between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to
efficiently load the same components into multiple pipelines.
Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section
[here](../../using-diffusers/svd#reduce-memory-usage).
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines. Also, to know more about reducing the memory usage of this pipeline, refer to the ["Reduce memory usage"] section [here](../../using-diffusers/svd#reduce-memory-usage).
</Tip>
<Tip warning={true}>
Marigold pipelines were designed and tested with the scheduler embedded in the model checkpoint.
The optimal number of inference steps varies by scheduler, with no universal value that works best across all cases.
To accommodate this, the `num_inference_steps` parameter in the pipeline's `__call__` method defaults to `None` (see the
API reference).
Unless set explicitly, it inherits the value from the `default_denoising_steps` field in the checkpoint configuration
file (`model_index.json`).
This ensures high-quality predictions when invoking the pipeline with only the `image` argument.
Marigold pipelines were designed and tested only with `DDIMScheduler` and `LCMScheduler`.
Depending on the scheduler, the number of inference steps required to get reliable predictions varies, and there is no universal value that works best across schedulers.
Because of that, the default value of `num_inference_steps` in the `__call__` method of the pipeline is set to `None` (see the API reference).
Unless set explicitly, its value will be taken from the checkpoint configuration `model_index.json`.
This is done to ensure high-quality predictions when calling the pipeline with just the `image` argument.
</Tip>
See also Marigold [usage examples](../../using-diffusers/marigold_usage).
## Marigold Depth Prediction API
See also Marigold [usage examples](marigold_usage).
## MarigoldDepthPipeline
[[autodoc]] MarigoldDepthPipeline
- all
- __call__
## MarigoldNormalsPipeline
[[autodoc]] MarigoldNormalsPipeline
- all
- __call__
## MarigoldDepthOutput
[[autodoc]] pipelines.marigold.pipeline_marigold_depth.MarigoldDepthOutput
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_depth
## Marigold Normals Estimation API
[[autodoc]] MarigoldNormalsPipeline
- __call__
[[autodoc]] pipelines.marigold.pipeline_marigold_normals.MarigoldNormalsOutput
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_normals
## Marigold Intrinsic Image Decomposition API
[[autodoc]] MarigoldIntrinsicsPipeline
- __call__
[[autodoc]] pipelines.marigold.pipeline_marigold_intrinsics.MarigoldIntrinsicsOutput
[[autodoc]] pipelines.marigold.marigold_image_processing.MarigoldImageProcessor.visualize_intrinsics
## MarigoldNormalsOutput
[[autodoc]] pipelines.marigold.pipeline_marigold_normals.MarigoldNormalsOutput

View File

@@ -15,10 +15,6 @@
# Mochi 1 Preview
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
> [!TIP]
> Only a research preview of the model weights is available at the moment.
@@ -119,7 +115,7 @@ export_to_video(frames, "mochi.mp4", fps=30)
## Reproducing the results from the Genmo Mochi repo
The [Genmo Mochi implementation](https://github.com/genmoai/mochi/tree/main) uses different precision values for each stage in the inference process. The text encoder and VAE use `torch.float32`, while the DiT uses `torch.bfloat16` with the [attention kernel](https://pytorch.org/docs/stable/generated/torch.nn.attention.sdpa_kernel.html#torch.nn.attention.sdpa_kernel) set to `EFFICIENT_ATTENTION`. Diffusers pipelines currently do not support setting different `dtypes` for different stages of the pipeline. In order to run inference in the same way as the original implementation, please refer to the following example.
The [Genmo Mochi implementation](https://github.com/genmoai/mochi/tree/main) uses different precision values for each stage in the inference process. The text encoder and VAE use `torch.float32`, while the DiT uses `torch.bfloat16` with the [attention kernel](https://pytorch.org/docs/stable/generated/torch.nn.attention.sdpa_kernel.html#torch.nn.attention.sdpa_kernel) set to `EFFICIENT_ATTENTION`. Diffusers pipelines currently do not support setting different `dtypes` for different stages of the pipeline. In order to run inference in the same way as the the original implementation, please refer to the following example.
<Tip>
The original Mochi implementation zeros out empty prompts. However, enabling this option and placing the entire pipeline under autocast can lead to numerical overflows with the T5 text encoder.

View File

@@ -1,80 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# OmniGen
[OmniGen: Unified Image Generation](https://arxiv.org/pdf/2409.11340) from BAAI, by Shitao Xiao, Yueze Wang, Junjie Zhou, Huaying Yuan, Xingrun Xing, Ruiran Yan, Chaofan Li, Shuting Wang, Tiejun Huang, Zheng Liu.
The abstract from the paper is:
*The emergence of Large Language Models (LLMs) has unified language generation tasks and revolutionized human-machine interaction. However, in the realm of image generation, a unified model capable of handling various tasks within a single framework remains largely unexplored. In this work, we introduce OmniGen, a new diffusion model for unified image generation. OmniGen is characterized by the following features: 1) Unification: OmniGen not only demonstrates text-to-image generation capabilities but also inherently supports various downstream tasks, such as image editing, subject-driven generation, and visual conditional generation. 2) Simplicity: The architecture of OmniGen is highly simplified, eliminating the need for additional plugins. Moreover, compared to existing diffusion models, it is more user-friendly and can complete complex tasks end-to-end through instructions without the need for extra intermediate steps, greatly simplifying the image generation workflow. 3) Knowledge Transfer: Benefit from learning in a unified format, OmniGen effectively transfers knowledge across different tasks, manages unseen tasks and domains, and exhibits novel capabilities. We also explore the models reasoning capabilities and potential applications of the chain-of-thought mechanism. This work represents the first attempt at a general-purpose image generation model, and we will release our resources at https://github.com/VectorSpaceLab/OmniGen to foster future advancements.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers.md) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading.md#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [staoxiao](https://github.com/staoxiao). The original codebase can be found [here](https://github.com/VectorSpaceLab/OmniGen). The original weights can be found under [hf.co/shitao](https://huggingface.co/Shitao/OmniGen-v1).
## Inference
First, load the pipeline:
```python
import torch
from diffusers import OmniGenPipeline
pipe = OmniGenPipeline.from_pretrained("Shitao/OmniGen-v1-diffusers", torch_dtype=torch.bfloat16)
pipe.to("cuda")
```
For text-to-image, pass a text prompt. By default, OmniGen generates a 1024x1024 image.
You can try setting the `height` and `width` parameters to generate images with different size.
```python
prompt = "Realistic photo. A young woman sits on a sofa, holding a book and facing the camera. She wears delicate silver hoop earrings adorned with tiny, sparkling diamonds that catch the light, with her long chestnut hair cascading over her shoulders. Her eyes are focused and gentle, framed by long, dark lashes. She is dressed in a cozy cream sweater, which complements her warm, inviting smile. Behind her, there is a table with a cup of water in a sleek, minimalist blue mug. The background is a serene indoor setting with soft natural light filtering through a window, adorned with tasteful art and flowers, creating a cozy and peaceful ambiance. 4K, HD."
image = pipe(
prompt=prompt,
height=1024,
width=1024,
guidance_scale=3,
generator=torch.Generator(device="cpu").manual_seed(111),
).images[0]
image.save("output.png")
```
OmniGen supports multimodal inputs.
When the input includes an image, you need to add a placeholder `<img><|image_1|></img>` in the text prompt to represent the image.
It is recommended to enable `use_input_image_size_as_output` to keep the edited image the same size as the original image.
```python
prompt="<img><|image_1|></img> Remove the woman's earrings. Replace the mug with a clear glass filled with sparkling iced cola."
input_images=[load_image("https://raw.githubusercontent.com/VectorSpaceLab/OmniGen/main/imgs/docs_img/t2i_woman_with_book.png")]
image = pipe(
prompt=prompt,
input_images=input_images,
guidance_scale=2,
img_guidance_scale=1.6,
use_input_image_size_as_output=True,
generator=torch.Generator(device="cpu").manual_seed(222)).images[0]
image.save("output.png")
```
## OmniGenPipeline
[[autodoc]] OmniGenPipeline
- all
- __call__

View File

@@ -54,7 +54,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [DiT](dit) | text2image |
| [Flux](flux) | text2image |
| [Hunyuan-DiT](hunyuandit) | text2image |
| [I2VGen-XL](i2vgenxl) | image2video |
| [I2VGen-XL](i2vgenxl) | text2video |
| [InstructPix2Pix](pix2pix) | image editing |
| [Kandinsky 2.1](kandinsky) | text2image, image2image, inpainting, interpolation |
| [Kandinsky 2.2](kandinsky_v22) | text2image, image2image, inpainting |
@@ -65,7 +65,7 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [Latte](latte) | text2image |
| [LEDITS++](ledits_pp) | image editing |
| [Lumina-T2X](lumina) | text2image |
| [Marigold](marigold) | depth-estimation, normals-estimation, intrinsic-decomposition |
| [Marigold](marigold) | depth |
| [MultiDiffusion](panorama) | text2image |
| [MusicLDM](musicldm) | text2audio |
| [PAG](pag) | text2image |
@@ -89,7 +89,6 @@ The table below lists all the pipelines currently available in 🤗 Diffusers an
| [UniDiffuser](unidiffuser) | text2image, image2text, image variation, text variation, unconditional image generation, unconditional audio generation |
| [Value-guided planning](value_guided_sampling) | value guided sampling |
| [Wuerstchen](wuerstchen) | text2image |
| [VisualCloze](visualcloze) | text2image, image2image, subject driven generation, inpainting, style transfer, image restoration, image editing, [depth,normal,edge,pose]2image, [depth,normal,edge,pose]-estimation, virtual try-on, image relighting |
## DiffusionPipeline

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Perturbed-Attention Guidance
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Perturbed-Attention Guidance (PAG)](https://ku-cvlab.github.io/Perturbed-Attention-Guidance/) is a new diffusion sampling guidance that improves sample quality across both unconditional and conditional settings, achieving this without requiring further training or the integration of external modules.
PAG was introduced in [Self-Rectifying Diffusion Sampling with Perturbed-Attention Guidance](https://huggingface.co/papers/2403.17377) by Donghoon Ahn, Hyoungwon Cho, Jaewon Min, Wooseok Jang, Jungwoo Kim, SeonHwa Kim, Hyun Hee Park, Kyong Hwan Jin and Seungryong Kim.

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# MultiDiffusion
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[MultiDiffusion: Fusing Diffusion Paths for Controlled Image Generation](https://huggingface.co/papers/2302.08113) is by Omer Bar-Tal, Lior Yariv, Yaron Lipman, and Tali Dekel.
The abstract from the paper is:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Image-to-Video Generation with PIA (Personalized Image Animator)
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
## Overview
[PIA: Your Personalized Image Animator via Plug-and-Play Modules in Text-to-Image Models](https://arxiv.org/abs/2312.13964) by Yiming Zhang, Zhening Xing, Yanhong Zeng, Youqing Fang, Kai Chen

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# InstructPix2Pix
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[InstructPix2Pix: Learning to Follow Image Editing Instructions](https://huggingface.co/papers/2211.09800) is by Tim Brooks, Aleksander Holynski and Alexei A. Efros.
The abstract from the paper is:

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@@ -14,11 +14,6 @@
# SanaPipeline
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
[SANA: Efficient High-Resolution Image Synthesis with Linear Diffusion Transformers](https://huggingface.co/papers/2410.10629) from NVIDIA and MIT HAN Lab, by Enze Xie, Junsong Chen, Junyu Chen, Han Cai, Haotian Tang, Yujun Lin, Zhekai Zhang, Muyang Li, Ligeng Zhu, Yao Lu, Song Han.
The abstract from the paper is:
@@ -64,10 +59,10 @@ Refer to the [Quantization](../../quantization/overview) overview to learn more
```py
import torch
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, SanaTransformer2DModel, SanaPipeline
from transformers import BitsAndBytesConfig as BitsAndBytesConfig, AutoModel
from transformers import BitsAndBytesConfig as BitsAndBytesConfig, AutoModelForCausalLM
quant_config = BitsAndBytesConfig(load_in_8bit=True)
text_encoder_8bit = AutoModel.from_pretrained(
text_encoder_8bit = AutoModelForCausalLM.from_pretrained(
"Efficient-Large-Model/Sana_1600M_1024px_diffusers",
subfolder="text_encoder",
quantization_config=quant_config,

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@@ -1,100 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# SANA-Sprint
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[SANA-Sprint: One-Step Diffusion with Continuous-Time Consistency Distillation](https://huggingface.co/papers/2503.09641) from NVIDIA, MIT HAN Lab, and Hugging Face by Junsong Chen, Shuchen Xue, Yuyang Zhao, Jincheng Yu, Sayak Paul, Junyu Chen, Han Cai, Enze Xie, Song Han
The abstract from the paper is:
*This paper presents SANA-Sprint, an efficient diffusion model for ultra-fast text-to-image (T2I) generation. SANA-Sprint is built on a pre-trained foundation model and augmented with hybrid distillation, dramatically reducing inference steps from 20 to 1-4. We introduce three key innovations: (1) We propose a training-free approach that transforms a pre-trained flow-matching model for continuous-time consistency distillation (sCM), eliminating costly training from scratch and achieving high training efficiency. Our hybrid distillation strategy combines sCM with latent adversarial distillation (LADD): sCM ensures alignment with the teacher model, while LADD enhances single-step generation fidelity. (2) SANA-Sprint is a unified step-adaptive model that achieves high-quality generation in 1-4 steps, eliminating step-specific training and improving efficiency. (3) We integrate ControlNet with SANA-Sprint for real-time interactive image generation, enabling instant visual feedback for user interaction. SANA-Sprint establishes a new Pareto frontier in speed-quality tradeoffs, achieving state-of-the-art performance with 7.59 FID and 0.74 GenEval in only 1 step — outperforming FLUX-schnell (7.94 FID / 0.71 GenEval) while being 10× faster (0.1s vs 1.1s on H100). It also achieves 0.1s (T2I) and 0.25s (ControlNet) latency for 1024×1024 images on H100, and 0.31s (T2I) on an RTX 4090, showcasing its exceptional efficiency and potential for AI-powered consumer applications (AIPC). Code and pre-trained models will be open-sourced.*
<Tip>
Make sure to check out the Schedulers [guide](../../using-diffusers/schedulers) to learn how to explore the tradeoff between scheduler speed and quality, and see the [reuse components across pipelines](../../using-diffusers/loading#reuse-a-pipeline) section to learn how to efficiently load the same components into multiple pipelines.
</Tip>
This pipeline was contributed by [lawrence-cj](https://github.com/lawrence-cj), [shuchen Xue](https://github.com/scxue) and [Enze Xie](https://github.com/xieenze). The original codebase can be found [here](https://github.com/NVlabs/Sana). The original weights can be found under [hf.co/Efficient-Large-Model](https://huggingface.co/Efficient-Large-Model/).
Available models:
| Model | Recommended dtype |
|:-------------------------------------------------------------------------------------------------------------------------------------------:|:-----------------:|
| [`Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers) | `torch.bfloat16` |
| [`Efficient-Large-Model/Sana_Sprint_0.6B_1024px_diffusers`](https://huggingface.co/Efficient-Large-Model/Sana_Sprint_0.6B_1024px_diffusers) | `torch.bfloat16` |
Refer to [this](https://huggingface.co/collections/Efficient-Large-Model/sana-sprint-67d6810d65235085b3b17c76) collection for more information.
Note: The recommended dtype mentioned is for the transformer weights. The text encoder must stay in `torch.bfloat16` and VAE weights must stay in `torch.bfloat16` or `torch.float32` for the model to work correctly. Please refer to the inference example below to see how to load the model with the recommended dtype.
## Quantization
Quantization helps reduce the memory requirements of very large models by storing model weights in a lower precision data type. However, quantization may have varying impact on video quality depending on the video model.
Refer to the [Quantization](../../quantization/overview) overview to learn more about supported quantization backends and selecting a quantization backend that supports your use case. The example below demonstrates how to load a quantized [`SanaSprintPipeline`] for inference with bitsandbytes.
```py
import torch
from diffusers import BitsAndBytesConfig as DiffusersBitsAndBytesConfig, SanaTransformer2DModel, SanaSprintPipeline
from transformers import BitsAndBytesConfig as BitsAndBytesConfig, AutoModel
quant_config = BitsAndBytesConfig(load_in_8bit=True)
text_encoder_8bit = AutoModel.from_pretrained(
"Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers",
subfolder="text_encoder",
quantization_config=quant_config,
torch_dtype=torch.bfloat16,
)
quant_config = DiffusersBitsAndBytesConfig(load_in_8bit=True)
transformer_8bit = SanaTransformer2DModel.from_pretrained(
"Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers",
subfolder="transformer",
quantization_config=quant_config,
torch_dtype=torch.bfloat16,
)
pipeline = SanaSprintPipeline.from_pretrained(
"Efficient-Large-Model/Sana_Sprint_1.6B_1024px_diffusers",
text_encoder=text_encoder_8bit,
transformer=transformer_8bit,
torch_dtype=torch.bfloat16,
device_map="balanced",
)
prompt = "a tiny astronaut hatching from an egg on the moon"
image = pipeline(prompt).images[0]
image.save("sana.png")
```
## Setting `max_timesteps`
Users can tweak the `max_timesteps` value for experimenting with the visual quality of the generated outputs. The default `max_timesteps` value was obtained with an inference-time search process. For more details about it, check out the paper.
## SanaSprintPipeline
[[autodoc]] SanaSprintPipeline
- all
- __call__
## SanaPipelineOutput
[[autodoc]] pipelines.sana.pipeline_output.SanaPipelineOutput

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Depth-to-image
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model can also infer depth based on an image using [MiDaS](https://github.com/isl-org/MiDaS). This allows you to pass a text prompt and an initial image to condition the generation of new images as well as a `depth_map` to preserve the image structure.
<Tip>

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Image-to-image
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model can also be applied to image-to-image generation by passing a text prompt and an initial image to condition the generation of new images.
The [`StableDiffusionImg2ImgPipeline`] uses the diffusion-denoising mechanism proposed in [SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equations](https://huggingface.co/papers/2108.01073) by Chenlin Meng, Yutong He, Yang Song, Jiaming Song, Jiajun Wu, Jun-Yan Zhu, Stefano Ermon.

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Inpainting
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model can also be applied to inpainting which lets you edit specific parts of an image by providing a mask and a text prompt using Stable Diffusion.
## Tips

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text-to-(RGB, depth)
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
LDM3D was proposed in [LDM3D: Latent Diffusion Model for 3D](https://huggingface.co/papers/2305.10853) by Gabriela Ben Melech Stan, Diana Wofk, Scottie Fox, Alex Redden, Will Saxton, Jean Yu, Estelle Aflalo, Shao-Yen Tseng, Fabio Nonato, Matthias Muller, and Vasudev Lal. LDM3D generates an image and a depth map from a given text prompt unlike the existing text-to-image diffusion models such as [Stable Diffusion](./overview) which only generates an image. With almost the same number of parameters, LDM3D achieves to create a latent space that can compress both the RGB images and the depth maps.
Two checkpoints are available for use:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable Diffusion pipelines
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/) and [LAION](https://laion.ai/). Latent diffusion applies the diffusion process over a lower dimensional latent space to reduce memory and compute complexity. This specific type of diffusion model was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
Stable Diffusion is trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs.

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@@ -12,11 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable Diffusion 3
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
Stable Diffusion 3 (SD3) was proposed in [Scaling Rectified Flow Transformers for High-Resolution Image Synthesis](https://arxiv.org/pdf/2403.03206.pdf) by Patrick Esser, Sumith Kulal, Andreas Blattmann, Rahim Entezari, Jonas Muller, Harry Saini, Yam Levi, Dominik Lorenz, Axel Sauer, Frederic Boesel, Dustin Podell, Tim Dockhorn, Zion English, Kyle Lacey, Alex Goodwin, Yannik Marek, and Robin Rombach.
The abstract from the paper is:

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@@ -12,11 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable Diffusion XL
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
<img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22">
</div>
Stable Diffusion XL (SDXL) was proposed in [SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis](https://huggingface.co/papers/2307.01952) by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, and Robin Rombach.
The abstract from the paper is:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text-to-image
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion model was created by researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), [Runway](https://github.com/runwayml), and [LAION](https://laion.ai/). The [`StableDiffusionPipeline`] is capable of generating photorealistic images given any text input. It's trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight and can run on consumer GPUs. Latent diffusion is the research on top of which Stable Diffusion was built. It was proposed in [High-Resolution Image Synthesis with Latent Diffusion Models](https://huggingface.co/papers/2112.10752) by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
The abstract from the paper is:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Super-resolution
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The Stable Diffusion upscaler diffusion model was created by the researchers and engineers from [CompVis](https://github.com/CompVis), [Stability AI](https://stability.ai/), and [LAION](https://laion.ai/). It is used to enhance the resolution of input images by a factor of 4.
<Tip>

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Stable unCLIP
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
Stable unCLIP checkpoints are finetuned from [Stable Diffusion 2.1](./stable_diffusion/stable_diffusion_2) checkpoints to condition on CLIP image embeddings.
Stable unCLIP still conditions on text embeddings. Given the two separate conditionings, stable unCLIP can be used
for text guided image variation. When combined with an unCLIP prior, it can also be used for full text to image generation.

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@@ -18,10 +18,6 @@ specific language governing permissions and limitations under the License.
# Text-to-video
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[ModelScope Text-to-Video Technical Report](https://arxiv.org/abs/2308.06571) is by Jiuniu Wang, Hangjie Yuan, Dayou Chen, Yingya Zhang, Xiang Wang, Shiwei Zhang.
The abstract from the paper is:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Text2Video-Zero
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Text2Video-Zero: Text-to-Image Diffusion Models are Zero-Shot Video Generators](https://huggingface.co/papers/2303.13439) is by Levon Khachatryan, Andranik Movsisyan, Vahram Tadevosyan, Roberto Henschel, [Zhangyang Wang](https://www.ece.utexas.edu/people/faculty/atlas-wang), Shant Navasardyan, [Humphrey Shi](https://www.humphreyshi.com).
Text2Video-Zero enables zero-shot video generation using either:

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@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# UniDiffuser
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
The UniDiffuser model was proposed in [One Transformer Fits All Distributions in Multi-Modal Diffusion at Scale](https://huggingface.co/papers/2303.06555) by Fan Bao, Shen Nie, Kaiwen Xue, Chongxuan Li, Shi Pu, Yaole Wang, Gang Yue, Yue Cao, Hang Su, Jun Zhu.
The abstract from the paper is:

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@@ -1,300 +0,0 @@
<!--Copyright 2025 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License.
-->
# VisualCloze
[VisualCloze: A Universal Image Generation Framework via Visual In-Context Learning](https://arxiv.org/abs/2504.07960) is an innovative in-context learning based universal image generation framework that offers key capabilities:
1. Support for various in-domain tasks
2. Generalization to unseen tasks through in-context learning
3. Unify multiple tasks into one step and generate both target image and intermediate results
4. Support reverse-engineering conditions from target images
## Overview
The abstract from the paper is:
*Recent progress in diffusion models significantly advances various image generation tasks. However, the current mainstream approach remains focused on building task-specific models, which have limited efficiency when supporting a wide range of different needs. While universal models attempt to address this limitation, they face critical challenges, including generalizable task instruction, appropriate task distributions, and unified architectural design. To tackle these challenges, we propose VisualCloze, a universal image generation framework, which supports a wide range of in-domain tasks, generalization to unseen ones, unseen unification of multiple tasks, and reverse generation. Unlike existing methods that rely on language-based task instruction, leading to task ambiguity and weak generalization, we integrate visual in-context learning, allowing models to identify tasks from visual demonstrations. Meanwhile, the inherent sparsity of visual task distributions hampers the learning of transferable knowledge across tasks. To this end, we introduce Graph200K, a graph-structured dataset that establishes various interrelated tasks, enhancing task density and transferable knowledge. Furthermore, we uncover that our unified image generation formulation shared a consistent objective with image infilling, enabling us to leverage the strong generative priors of pre-trained infilling models without modifying the architectures. The codes, dataset, and models are available at https://visualcloze.github.io.*
## Inference
### Model loading
VisualCloze is a two-stage cascade pipeline, containing `VisualClozeGenerationPipeline` and `VisualClozeUpsamplingPipeline`.
- In `VisualClozeGenerationPipeline`, each image is downsampled before concatenating images into a grid layout, avoiding excessively high resolutions. VisualCloze releases two models suitable for diffusers, i.e., [VisualClozePipeline-384](https://huggingface.co/VisualCloze/VisualClozePipeline-384) and [VisualClozePipeline-512](https://huggingface.co/VisualCloze/VisualClozePipeline-384), which downsample images to resolutions of 384 and 512, respectively.
- `VisualClozeUpsamplingPipeline` uses [SDEdit](https://arxiv.org/abs/2108.01073) to enable high-resolution image synthesis.
The `VisualClozePipeline` integrates both stages to support convenient end-to-end sampling, while also allowing users to utilize each pipeline independently as needed.
### Input Specifications
#### Task and Content Prompts
- Task prompt: Required to describe the generation task intention
- Content prompt: Optional description or caption of the target image
- When content prompt is not needed, pass `None`
- For batch inference, pass `List[str|None]`
#### Image Input Format
- Format: `List[List[Image|None]]`
- Structure:
- All rows except the last represent in-context examples
- Last row represents the current query (target image set to `None`)
- For batch inference, pass `List[List[List[Image|None]]]`
#### Resolution Control
- Default behavior:
- Initial generation in the first stage: area of ${pipe.resolution}^2$
- Upsampling in the second stage: 3x factor
- Custom resolution: Adjust using `upsampling_height` and `upsampling_width` parameters
### Examples
For comprehensive examples covering a wide range of tasks, please refer to the [Online Demo](https://huggingface.co/spaces/VisualCloze/VisualCloze) and [GitHub Repository](https://github.com/lzyhha/VisualCloze). Below are simple examples for three cases: mask-to-image conversion, edge detection, and subject-driven generation.
#### Example for mask2image
```python
import torch
from diffusers import VisualClozePipeline
from diffusers.utils import load_image
pipe = VisualClozePipeline.from_pretrained("VisualCloze/VisualClozePipeline-384", resolution=384, torch_dtype=torch.bfloat16)
pipe.to("cuda")
# Load in-context images (make sure the paths are correct and accessible)
image_paths = [
# in-context examples
[
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_mask2image_incontext-example-1_mask.jpg'),
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_mask2image_incontext-example-1_image.jpg'),
],
# query with the target image
[
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_mask2image_query_mask.jpg'),
None, # No image needed for the target image
],
]
# Task and content prompt
task_prompt = "In each row, a logical task is demonstrated to achieve [IMAGE2] an aesthetically pleasing photograph based on [IMAGE1] sam 2-generated masks with rich color coding."
content_prompt = """Majestic photo of a golden eagle perched on a rocky outcrop in a mountainous landscape.
The eagle is positioned in the right foreground, facing left, with its sharp beak and keen eyes prominently visible.
Its plumage is a mix of dark brown and golden hues, with intricate feather details.
The background features a soft-focus view of snow-capped mountains under a cloudy sky, creating a serene and grandiose atmosphere.
The foreground includes rugged rocks and patches of green moss. Photorealistic, medium depth of field,
soft natural lighting, cool color palette, high contrast, sharp focus on the eagle, blurred background,
tranquil, majestic, wildlife photography."""
# Run the pipeline
image_result = pipe(
task_prompt=task_prompt,
content_prompt=content_prompt,
image=image_paths,
upsampling_width=1344,
upsampling_height=768,
upsampling_strength=0.4,
guidance_scale=30,
num_inference_steps=30,
max_sequence_length=512,
generator=torch.Generator("cpu").manual_seed(0)
).images[0][0]
# Save the resulting image
image_result.save("visualcloze.png")
```
#### Example for edge-detection
```python
import torch
from diffusers import VisualClozePipeline
from diffusers.utils import load_image
pipe = VisualClozePipeline.from_pretrained("VisualCloze/VisualClozePipeline-384", resolution=384, torch_dtype=torch.bfloat16)
pipe.to("cuda")
# Load in-context images (make sure the paths are correct and accessible)
image_paths = [
# in-context examples
[
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_edgedetection_incontext-example-1_image.jpg'),
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_edgedetection_incontext-example-1_edge.jpg'),
],
[
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_edgedetection_incontext-example-2_image.jpg'),
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_edgedetection_incontext-example-2_edge.jpg'),
],
# query with the target image
[
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_edgedetection_query_image.jpg'),
None, # No image needed for the target image
],
]
# Task and content prompt
task_prompt = "Each row illustrates a pathway from [IMAGE1] a sharp and beautifully composed photograph to [IMAGE2] edge map with natural well-connected outlines using a clear logical task."
content_prompt = ""
# Run the pipeline
image_result = pipe(
task_prompt=task_prompt,
content_prompt=content_prompt,
image=image_paths,
upsampling_width=864,
upsampling_height=1152,
upsampling_strength=0.4,
guidance_scale=30,
num_inference_steps=30,
max_sequence_length=512,
generator=torch.Generator("cpu").manual_seed(0)
).images[0][0]
# Save the resulting image
image_result.save("visualcloze.png")
```
#### Example for subject-driven generation
```python
import torch
from diffusers import VisualClozePipeline
from diffusers.utils import load_image
pipe = VisualClozePipeline.from_pretrained("VisualCloze/VisualClozePipeline-384", resolution=384, torch_dtype=torch.bfloat16)
pipe.to("cuda")
# Load in-context images (make sure the paths are correct and accessible)
image_paths = [
# in-context examples
[
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_subjectdriven_incontext-example-1_reference.jpg'),
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_subjectdriven_incontext-example-1_depth.jpg'),
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_subjectdriven_incontext-example-1_image.jpg'),
],
[
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_subjectdriven_incontext-example-2_reference.jpg'),
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_subjectdriven_incontext-example-2_depth.jpg'),
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_subjectdriven_incontext-example-2_image.jpg'),
],
# query with the target image
[
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_subjectdriven_query_reference.jpg'),
load_image('https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_subjectdriven_query_depth.jpg'),
None, # No image needed for the target image
],
]
# Task and content prompt
task_prompt = """Each row describes a process that begins with [IMAGE1] an image containing the key object,
[IMAGE2] depth map revealing gray-toned spatial layers and results in
[IMAGE3] an image with artistic qualitya high-quality image with exceptional detail."""
content_prompt = """A vintage porcelain collector's item. Beneath a blossoming cherry tree in early spring,
this treasure is photographed up close, with soft pink petals drifting through the air and vibrant blossoms framing the scene."""
# Run the pipeline
image_result = pipe(
task_prompt=task_prompt,
content_prompt=content_prompt,
image=image_paths,
upsampling_width=1024,
upsampling_height=1024,
upsampling_strength=0.2,
guidance_scale=30,
num_inference_steps=30,
max_sequence_length=512,
generator=torch.Generator("cpu").manual_seed(0)
).images[0][0]
# Save the resulting image
image_result.save("visualcloze.png")
```
#### Utilize each pipeline independently
```python
import torch
from diffusers import VisualClozeGenerationPipeline, FluxFillPipeline as VisualClozeUpsamplingPipeline
from diffusers.utils import load_image
from PIL import Image
pipe = VisualClozeGenerationPipeline.from_pretrained(
"VisualCloze/VisualClozePipeline-384", resolution=384, torch_dtype=torch.bfloat16
)
pipe.to("cuda")
image_paths = [
# in-context examples
[
load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_mask2image_incontext-example-1_mask.jpg"
),
load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_mask2image_incontext-example-1_image.jpg"
),
],
# query with the target image
[
load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/visualcloze/visualcloze_mask2image_query_mask.jpg"
),
None, # No image needed for the target image
],
]
task_prompt = "In each row, a logical task is demonstrated to achieve [IMAGE2] an aesthetically pleasing photograph based on [IMAGE1] sam 2-generated masks with rich color coding."
content_prompt = "Majestic photo of a golden eagle perched on a rocky outcrop in a mountainous landscape. The eagle is positioned in the right foreground, facing left, with its sharp beak and keen eyes prominently visible. Its plumage is a mix of dark brown and golden hues, with intricate feather details. The background features a soft-focus view of snow-capped mountains under a cloudy sky, creating a serene and grandiose atmosphere. The foreground includes rugged rocks and patches of green moss. Photorealistic, medium depth of field, soft natural lighting, cool color palette, high contrast, sharp focus on the eagle, blurred background, tranquil, majestic, wildlife photography."
# Stage 1: Generate initial image
image = pipe(
task_prompt=task_prompt,
content_prompt=content_prompt,
image=image_paths,
guidance_scale=30,
num_inference_steps=30,
max_sequence_length=512,
generator=torch.Generator("cpu").manual_seed(0),
).images[0][0]
# Stage 2 (optional): Upsample the generated image
pipe_upsample = VisualClozeUpsamplingPipeline.from_pipe(pipe)
pipe_upsample.to("cuda")
mask_image = Image.new("RGB", image.size, (255, 255, 255))
image = pipe_upsample(
image=image,
mask_image=mask_image,
prompt=content_prompt,
width=1344,
height=768,
strength=0.4,
guidance_scale=30,
num_inference_steps=30,
max_sequence_length=512,
generator=torch.Generator("cpu").manual_seed(0),
).images[0]
image.save("visualcloze.png")
```
## VisualClozePipeline
[[autodoc]] VisualClozePipeline
- all
- __call__
## VisualClozeGenerationPipeline
[[autodoc]] VisualClozeGenerationPipeline
- all
- __call__

View File

@@ -1,519 +0,0 @@
<!-- Copyright 2024 The HuggingFace Team. All rights reserved.
#
# Licensed under the Apache License, Version 2.0 (the "License");
# you may not use this file except in compliance with the License.
# You may obtain a copy of the License at
#
# http://www.apache.org/licenses/LICENSE-2.0
#
# Unless required by applicable law or agreed to in writing, software
# distributed under the License is distributed on an "AS IS" BASIS,
# WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied.
# See the License for the specific language governing permissions and
# limitations under the License. -->
# Wan
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
[Wan 2.1](https://github.com/Wan-Video/Wan2.1) by the Alibaba Wan Team.
<!-- TODO(aryan): update abstract once paper is out -->
## Generating Videos with Wan 2.1
We will first need to install some additional dependencies.
```shell
pip install -u ftfy imageio-ffmpeg imageio
```
### Text to Video Generation
The following example requires 11GB VRAM to run and uses the smaller `Wan-AI/Wan2.1-T2V-1.3B-Diffusers` model. You can switch it out
for the larger `Wan2.1-I2V-14B-720P-Diffusers` or `Wan-AI/Wan2.1-I2V-14B-480P-Diffusers` if you have at least 35GB VRAM available.
```python
from diffusers import WanPipeline
from diffusers.utils import export_to_video
# Available models: Wan-AI/Wan2.1-I2V-14B-720P-Diffusers or Wan-AI/Wan2.1-I2V-14B-480P-Diffusers
model_id = "Wan-AI/Wan2.1-T2V-1.3B-Diffusers"
pipe = WanPipeline.from_pretrained(model_id, torch_dtype=torch.bfloat16)
pipe.enable_model_cpu_offload()
prompt = "A cat and a dog baking a cake together in a kitchen. The cat is carefully measuring flour, while the dog is stirring the batter with a wooden spoon. The kitchen is cozy, with sunlight streaming through the window."
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
frames = pipe(prompt=prompt, negative_prompt=negative_prompt, num_frames=num_frames).frames[0]
export_to_video(frames, "wan-t2v.mp4", fps=16)
```
<Tip>
You can improve the quality of the generated video by running the decoding step in full precision.
</Tip>
```python
from diffusers import WanPipeline, AutoencoderKLWan
from diffusers.utils import export_to_video
model_id = "Wan-AI/Wan2.1-T2V-1.3B-Diffusers"
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
pipe = WanPipeline.from_pretrained(model_id, vae=vae, torch_dtype=torch.bfloat16)
# replace this with pipe.to("cuda") if you have sufficient VRAM
pipe.enable_model_cpu_offload()
prompt = "A cat and a dog baking a cake together in a kitchen. The cat is carefully measuring flour, while the dog is stirring the batter with a wooden spoon. The kitchen is cozy, with sunlight streaming through the window."
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
frames = pipe(prompt=prompt, num_frames=num_frames).frames[0]
export_to_video(frames, "wan-t2v.mp4", fps=16)
```
### Image to Video Generation
The Image to Video pipeline requires loading the `AutoencoderKLWan` and the `CLIPVisionModel` components in full precision. The following example will need at least
35GB of VRAM to run.
```python
import torch
import numpy as np
from diffusers import AutoencoderKLWan, WanImageToVideoPipeline
from diffusers.utils import export_to_video, load_image
from transformers import CLIPVisionModel
# Available models: Wan-AI/Wan2.1-I2V-14B-480P-Diffusers, Wan-AI/Wan2.1-I2V-14B-720P-Diffusers
model_id = "Wan-AI/Wan2.1-I2V-14B-480P-Diffusers"
image_encoder = CLIPVisionModel.from_pretrained(
model_id, subfolder="image_encoder", torch_dtype=torch.float32
)
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
pipe = WanImageToVideoPipeline.from_pretrained(
model_id, vae=vae, image_encoder=image_encoder, torch_dtype=torch.bfloat16
)
# replace this with pipe.to("cuda") if you have sufficient VRAM
pipe.enable_model_cpu_offload()
image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"
)
max_area = 480 * 832
aspect_ratio = image.height / image.width
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
image = image.resize((width, height))
prompt = (
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
)
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
output = pipe(
image=image,
prompt=prompt,
negative_prompt=negative_prompt,
height=height,
width=width,
num_frames=num_frames,
guidance_scale=5.0,
).frames[0]
export_to_video(output, "wan-i2v.mp4", fps=16)
```
### First and Last Frame Interpolation
```python
import numpy as np
import torch
import torchvision.transforms.functional as TF
from diffusers import AutoencoderKLWan, WanImageToVideoPipeline
from diffusers.utils import export_to_video, load_image
from transformers import CLIPVisionModel
model_id = "Wan-AI/Wan2.1-FLF2V-14B-720P-diffusers"
image_encoder = CLIPVisionModel.from_pretrained(model_id, subfolder="image_encoder", torch_dtype=torch.float32)
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
pipe = WanImageToVideoPipeline.from_pretrained(
model_id, vae=vae, image_encoder=image_encoder, torch_dtype=torch.bfloat16
)
pipe.to("cuda")
first_frame = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flf2v_input_first_frame.png")
last_frame = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/flf2v_input_last_frame.png")
def aspect_ratio_resize(image, pipe, max_area=720 * 1280):
aspect_ratio = image.height / image.width
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
image = image.resize((width, height))
return image, height, width
def center_crop_resize(image, height, width):
# Calculate resize ratio to match first frame dimensions
resize_ratio = max(width / image.width, height / image.height)
# Resize the image
width = round(image.width * resize_ratio)
height = round(image.height * resize_ratio)
size = [width, height]
image = TF.center_crop(image, size)
return image, height, width
first_frame, height, width = aspect_ratio_resize(first_frame, pipe)
if last_frame.size != first_frame.size:
last_frame, _, _ = center_crop_resize(last_frame, height, width)
prompt = "CG animation style, a small blue bird takes off from the ground, flapping its wings. The bird's feathers are delicate, with a unique pattern on its chest. The background shows a blue sky with white clouds under bright sunshine. The camera follows the bird upward, capturing its flight and the vastness of the sky from a close-up, low-angle perspective."
output = pipe(
image=first_frame, last_image=last_frame, prompt=prompt, height=height, width=width, guidance_scale=5.5
).frames[0]
export_to_video(output, "output.mp4", fps=16)
```
### Video to Video Generation
```python
import torch
from diffusers.utils import load_video, export_to_video
from diffusers import AutoencoderKLWan, WanVideoToVideoPipeline, UniPCMultistepScheduler
# Available models: Wan-AI/Wan2.1-T2V-14B-Diffusers, Wan-AI/Wan2.1-T2V-1.3B-Diffusers
model_id = "Wan-AI/Wan2.1-T2V-1.3B-Diffusers"
vae = AutoencoderKLWan.from_pretrained(
model_id, subfolder="vae", torch_dtype=torch.float32
)
pipe = WanVideoToVideoPipeline.from_pretrained(
model_id, vae=vae, torch_dtype=torch.bfloat16
)
flow_shift = 3.0 # 5.0 for 720P, 3.0 for 480P
pipe.scheduler = UniPCMultistepScheduler.from_config(
pipe.scheduler.config, flow_shift=flow_shift
)
# change to pipe.to("cuda") if you have sufficient VRAM
pipe.enable_model_cpu_offload()
prompt = "A robot standing on a mountain top. The sun is setting in the background"
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
video = load_video(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/hiker.mp4"
)
output = pipe(
video=video,
prompt=prompt,
negative_prompt=negative_prompt,
height=480,
width=512,
guidance_scale=7.0,
strength=0.7,
).frames[0]
export_to_video(output, "wan-v2v.mp4", fps=16)
```
## Memory Optimizations for Wan 2.1
Base inference with the large 14B Wan 2.1 models can take up to 35GB of VRAM when generating videos at 720p resolution. We'll outline a few memory optimizations we can apply to reduce the VRAM required to run the model.
We'll use `Wan-AI/Wan2.1-I2V-14B-720P-Diffusers` model in these examples to demonstrate the memory savings, but the techniques are applicable to all model checkpoints.
### Group Offloading the Transformer and UMT5 Text Encoder
Find more information about group offloading [here](../optimization/memory.md)
#### Block Level Group Offloading
We can reduce our VRAM requirements by applying group offloading to the larger model components of the pipeline; the `WanTransformer3DModel` and `UMT5EncoderModel`. Group offloading will break up the individual modules of a model and offload/onload them onto your GPU as needed during inference. In this example, we'll apply `block_level` offloading, which will group the modules in a model into blocks of size `num_blocks_per_group` and offload/onload them to GPU. Moving to between CPU and GPU does add latency to the inference process. You can trade off between latency and memory savings by increasing or decreasing the `num_blocks_per_group`.
The following example will now only require 14GB of VRAM to run, but will take approximately 30 minutes to generate a video.
```python
import torch
import numpy as np
from diffusers import AutoencoderKLWan, WanTransformer3DModel, WanImageToVideoPipeline
from diffusers.hooks.group_offloading import apply_group_offloading
from diffusers.utils import export_to_video, load_image
from transformers import UMT5EncoderModel, CLIPVisionModel
# Available models: Wan-AI/Wan2.1-I2V-14B-480P-Diffusers, Wan-AI/Wan2.1-I2V-14B-720P-Diffusers
model_id = "Wan-AI/Wan2.1-I2V-14B-720P-Diffusers"
image_encoder = CLIPVisionModel.from_pretrained(
model_id, subfolder="image_encoder", torch_dtype=torch.float32
)
text_encoder = UMT5EncoderModel.from_pretrained(model_id, subfolder="text_encoder", torch_dtype=torch.bfloat16)
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
transformer = WanTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
onload_device = torch.device("cuda")
offload_device = torch.device("cpu")
apply_group_offloading(text_encoder,
onload_device=onload_device,
offload_device=offload_device,
offload_type="block_level",
num_blocks_per_group=4
)
transformer.enable_group_offload(
onload_device=onload_device,
offload_device=offload_device,
offload_type="block_level",
num_blocks_per_group=4,
)
pipe = WanImageToVideoPipeline.from_pretrained(
model_id,
vae=vae,
transformer=transformer,
text_encoder=text_encoder,
image_encoder=image_encoder,
torch_dtype=torch.bfloat16
)
# Since we've offloaded the larger models already, we can move the rest of the model components to GPU
pipe.to("cuda")
image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"
)
max_area = 720 * 832
aspect_ratio = image.height / image.width
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
image = image.resize((width, height))
prompt = (
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
)
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
output = pipe(
image=image,
prompt=prompt,
negative_prompt=negative_prompt,
height=height,
width=width,
num_frames=num_frames,
guidance_scale=5.0,
).frames[0]
export_to_video(output, "wan-i2v.mp4", fps=16)
```
#### Block Level Group Offloading with CUDA Streams
We can speed up group offloading inference, by enabling the use of [CUDA streams](https://pytorch.org/docs/stable/generated/torch.cuda.Stream.html). However, using CUDA streams requires moving the model parameters into pinned memory. This allocation is handled by Pytorch under the hood, and can result in a significant spike in CPU RAM usage. Please consider this option if your CPU RAM is atleast 2X the size of the model you are group offloading.
In the following example we will use CUDA streams when group offloading the `WanTransformer3DModel`. When testing on an A100, this example will require 14GB of VRAM, 52GB of CPU RAM, but will generate a video in approximately 9 minutes.
```python
import torch
import numpy as np
from diffusers import AutoencoderKLWan, WanTransformer3DModel, WanImageToVideoPipeline
from diffusers.hooks.group_offloading import apply_group_offloading
from diffusers.utils import export_to_video, load_image
from transformers import UMT5EncoderModel, CLIPVisionModel
# Available models: Wan-AI/Wan2.1-I2V-14B-480P-Diffusers, Wan-AI/Wan2.1-I2V-14B-720P-Diffusers
model_id = "Wan-AI/Wan2.1-I2V-14B-720P-Diffusers"
image_encoder = CLIPVisionModel.from_pretrained(
model_id, subfolder="image_encoder", torch_dtype=torch.float32
)
text_encoder = UMT5EncoderModel.from_pretrained(model_id, subfolder="text_encoder", torch_dtype=torch.bfloat16)
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
transformer = WanTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
onload_device = torch.device("cuda")
offload_device = torch.device("cpu")
apply_group_offloading(text_encoder,
onload_device=onload_device,
offload_device=offload_device,
offload_type="block_level",
num_blocks_per_group=4
)
transformer.enable_group_offload(
onload_device=onload_device,
offload_device=offload_device,
offload_type="leaf_level",
use_stream=True
)
pipe = WanImageToVideoPipeline.from_pretrained(
model_id,
vae=vae,
transformer=transformer,
text_encoder=text_encoder,
image_encoder=image_encoder,
torch_dtype=torch.bfloat16
)
# Since we've offloaded the larger models already, we can move the rest of the model components to GPU
pipe.to("cuda")
image = load_image(
"https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"
)
max_area = 720 * 832
aspect_ratio = image.height / image.width
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
image = image.resize((width, height))
prompt = (
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
)
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
output = pipe(
image=image,
prompt=prompt,
negative_prompt=negative_prompt,
height=height,
width=width,
num_frames=num_frames,
guidance_scale=5.0,
).frames[0]
export_to_video(output, "wan-i2v.mp4", fps=16)
```
### Applying Layerwise Casting to the Transformer
Find more information about layerwise casting [here](../optimization/memory.md)
In this example, we will model offloading with layerwise casting. Layerwise casting will downcast each layer's weights to `torch.float8_e4m3fn`, temporarily upcast to `torch.bfloat16` during the forward pass of the layer, then revert to `torch.float8_e4m3fn` afterward. This approach reduces memory requirements by approximately 50% while introducing a minor quality reduction in the generated video due to the precision trade-off.
This example will require 20GB of VRAM.
```python
import torch
import numpy as np
from diffusers import AutoencoderKLWan, WanTransformer3DModel, WanImageToVideoPipeline
from diffusers.hooks.group_offloading import apply_group_offloading
from diffusers.utils import export_to_video, load_image
from transformers import UMT5EncoderModel, CLIPVisionModel
model_id = "Wan-AI/Wan2.1-I2V-14B-720P-Diffusers"
image_encoder = CLIPVisionModel.from_pretrained(
model_id, subfolder="image_encoder", torch_dtype=torch.float32
)
text_encoder = UMT5EncoderModel.from_pretrained(model_id, subfolder="text_encoder", torch_dtype=torch.bfloat16)
vae = AutoencoderKLWan.from_pretrained(model_id, subfolder="vae", torch_dtype=torch.float32)
transformer = WanTransformer3DModel.from_pretrained(model_id, subfolder="transformer", torch_dtype=torch.bfloat16)
transformer.enable_layerwise_casting(storage_dtype=torch.float8_e4m3fn, compute_dtype=torch.bfloat16)
pipe = WanImageToVideoPipeline.from_pretrained(
model_id,
vae=vae,
transformer=transformer,
text_encoder=text_encoder,
image_encoder=image_encoder,
torch_dtype=torch.bfloat16
)
pipe.enable_model_cpu_offload()
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg")
max_area = 720 * 832
aspect_ratio = image.height / image.width
mod_value = pipe.vae_scale_factor_spatial * pipe.transformer.config.patch_size[1]
height = round(np.sqrt(max_area * aspect_ratio)) // mod_value * mod_value
width = round(np.sqrt(max_area / aspect_ratio)) // mod_value * mod_value
image = image.resize((width, height))
prompt = (
"An astronaut hatching from an egg, on the surface of the moon, the darkness and depth of space realised in "
"the background. High quality, ultrarealistic detail and breath-taking movie-like camera shot."
)
negative_prompt = "Bright tones, overexposed, static, blurred details, subtitles, style, works, paintings, images, static, overall gray, worst quality, low quality, JPEG compression residue, ugly, incomplete, extra fingers, poorly drawn hands, poorly drawn faces, deformed, disfigured, misshapen limbs, fused fingers, still picture, messy background, three legs, many people in the background, walking backwards"
num_frames = 33
output = pipe(
image=image,
prompt=prompt,
negative_prompt=negative_prompt,
height=height,
width=width,
num_frames=num_frames,
num_inference_steps=50,
guidance_scale=5.0,
).frames[0]
export_to_video(output, "wan-i2v.mp4", fps=16)
```
## Using a Custom Scheduler
Wan can be used with many different schedulers, each with their own benefits regarding speed and generation quality. By default, Wan uses the `UniPCMultistepScheduler(prediction_type="flow_prediction", use_flow_sigmas=True, flow_shift=3.0)` scheduler. You can use a different scheduler as follows:
```python
from diffusers import FlowMatchEulerDiscreteScheduler, UniPCMultistepScheduler, WanPipeline
scheduler_a = FlowMatchEulerDiscreteScheduler(shift=5.0)
scheduler_b = UniPCMultistepScheduler(prediction_type="flow_prediction", use_flow_sigmas=True, flow_shift=4.0)
pipe = WanPipeline.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", scheduler=<CUSTOM_SCHEDULER_HERE>)
# or,
pipe.scheduler = <CUSTOM_SCHEDULER_HERE>
```
## Using Single File Loading with Wan 2.1
The `WanTransformer3DModel` and `AutoencoderKLWan` models support loading checkpoints in their original format via the `from_single_file` loading
method.
```python
import torch
from diffusers import WanPipeline, WanTransformer3DModel
ckpt_path = "https://huggingface.co/Comfy-Org/Wan_2.1_ComfyUI_repackaged/blob/main/split_files/diffusion_models/wan2.1_t2v_1.3B_bf16.safetensors"
transformer = WanTransformer3DModel.from_single_file(ckpt_path, torch_dtype=torch.bfloat16)
pipe = WanPipeline.from_pretrained("Wan-AI/Wan2.1-T2V-1.3B-Diffusers", transformer=transformer)
```
## Recommendations for Inference
- Keep `AutencoderKLWan` in `torch.float32` for better decoding quality.
- `num_frames` should satisfy the following constraint: `(num_frames - 1) % 4 == 0`
- For smaller resolution videos, try lower values of `shift` (between `2.0` to `5.0`) in the [Scheduler](https://huggingface.co/docs/diffusers/main/en/api/schedulers/flow_match_euler_discrete#diffusers.FlowMatchEulerDiscreteScheduler.shift). For larger resolution videos, try higher values (between `7.0` and `12.0`). The default value is `3.0` for Wan.
## WanPipeline
[[autodoc]] WanPipeline
- all
- __call__
## WanImageToVideoPipeline
[[autodoc]] WanImageToVideoPipeline
- all
- __call__
## WanPipelineOutput
[[autodoc]] pipelines.wan.pipeline_output.WanPipelineOutput

View File

@@ -12,10 +12,6 @@ specific language governing permissions and limitations under the License.
# Würstchen
<div class="flex flex-wrap space-x-1">
<img alt="LoRA" src="https://img.shields.io/badge/LoRA-d8b4fe?style=flat"/>
</div>
<img src="https://github.com/dome272/Wuerstchen/assets/61938694/0617c863-165a-43ee-9303-2a17299a0cf9">
[Wuerstchen: An Efficient Architecture for Large-Scale Text-to-Image Diffusion Models](https://huggingface.co/papers/2306.00637) is by Pablo Pernias, Dominic Rampas, Mats L. Richter and Christopher Pal and Marc Aubreville.

View File

@@ -13,7 +13,9 @@ specific language governing permissions and limitations under the License.
# Quantization
Quantization techniques reduce memory and computational costs by representing weights and activations with lower-precision data types like 8-bit integers (int8). This enables loading larger models you normally wouldn't be able to fit into memory, and speeding up inference.
Quantization techniques reduce memory and computational costs by representing weights and activations with lower-precision data types like 8-bit integers (int8). This enables loading larger models you normally wouldn't be able to fit into memory, and speeding up inference. Diffusers supports 8-bit and 4-bit quantization with [bitsandbytes](https://huggingface.co/docs/bitsandbytes/en/index).
Quantization techniques that aren't supported in Transformers can be added with the [`DiffusersQuantizer`] class.
<Tip>
@@ -21,9 +23,6 @@ Learn how to quantize models in the [Quantization](../quantization/overview) gui
</Tip>
## PipelineQuantizationConfig
[[autodoc]] quantizers.PipelineQuantizationConfig
## BitsAndBytesConfig
@@ -32,11 +31,6 @@ Learn how to quantize models in the [Quantization](../quantization/overview) gui
## GGUFQuantizationConfig
[[autodoc]] GGUFQuantizationConfig
## QuantoConfig
[[autodoc]] QuantoConfig
## TorchAoConfig
[[autodoc]] TorchAoConfig

View File

@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# CogVideoXDDIMScheduler
`CogVideoXDDIMScheduler` is based on [Denoising Diffusion Implicit Models](https://huggingface.co/papers/2010.02502), specifically for CogVideoX models.
## CogVideoXDDIMScheduler
[[autodoc]] CogVideoXDDIMScheduler

View File

@@ -1,19 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# CogVideoXDPMScheduler
`CogVideoXDPMScheduler` is based on [DPM-Solver: A Fast ODE Solver for Diffusion Probabilistic Model Sampling in Around 10 Steps](https://huggingface.co/papers/2206.00927) and [DPM-Solver++: Fast Solver for Guided Sampling of Diffusion Probabilistic Models](https://huggingface.co/papers/2211.01095), specifically for CogVideoX models.
## CogVideoXDPMScheduler
[[autodoc]] CogVideoXDPMScheduler

View File

@@ -41,11 +41,3 @@ Utility and helper functions for working with 🤗 Diffusers.
## randn_tensor
[[autodoc]] utils.torch_utils.randn_tensor
## apply_layerwise_casting
[[autodoc]] hooks.layerwise_casting.apply_layerwise_casting
## apply_group_offloading
[[autodoc]] hooks.group_offloading.apply_group_offloading

View File

@@ -83,8 +83,4 @@ Happy exploring, and thank you for being part of the Diffusers community!
<td><a href="https://github.com/suzukimain/auto_diffusers"> Model Search </a></td>
<td>Search models on Civitai and Hugging Face</td>
</tr>
<tr style="border-top: 2px solid black">
<td><a href="https://github.com/beinsezii/skrample"> Skrample </a></td>
<td>Fully modular scheduler functions with 1st class diffusers integration.</td>
</tr>
</table>

View File

@@ -16,11 +16,6 @@ specific language governing permissions and limitations under the License.
<img src="https://colab.research.google.com/assets/colab-badge.svg" alt="Open In Colab"/>
</a>
> [!TIP]
> This document has now grown outdated given the emergence of existing evaluation frameworks for diffusion models for image generation. Please check
> out works like [HEIM](https://crfm.stanford.edu/helm/heim/latest/), [T2I-Compbench](https://arxiv.org/abs/2307.06350),
> [GenEval](https://arxiv.org/abs/2310.11513).
Evaluation of generative models like [Stable Diffusion](https://huggingface.co/docs/diffusers/stable_diffusion) is subjective in nature. But as practitioners and researchers, we often have to make careful choices amongst many different possibilities. So, when working with different generative models (like GANs, Diffusion, etc.), how do we choose one over the other?
Qualitative evaluation of such models can be error-prone and might incorrectly influence a decision.

View File

@@ -1,9 +0,0 @@
# Hybrid Inference API Reference
## Remote Decode
[[autodoc]] utils.remote_utils.remote_decode
## Remote Encode
[[autodoc]] utils.remote_utils.remote_encode

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@@ -1,60 +0,0 @@
<!--Copyright 2024 The HuggingFace Team. All rights reserved.
Licensed under the Apache License, Version 2.0 (the "License"); you may not use this file except in compliance with
the License. You may obtain a copy of the License at
http://www.apache.org/licenses/LICENSE-2.0
Unless required by applicable law or agreed to in writing, software distributed under the License is distributed on
an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. See the License for the
specific language governing permissions and limitations under the License.
-->
# Hybrid Inference
**Empowering local AI builders with Hybrid Inference**
> [!TIP]
> Hybrid Inference is an [experimental feature](https://huggingface.co/blog/remote_vae).
> Feedback can be provided [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
## Why use Hybrid Inference?
Hybrid Inference offers a fast and simple way to offload local generation requirements.
- 🚀 **Reduced Requirements:** Access powerful models without expensive hardware.
- 💎 **Without Compromise:** Achieve the highest quality without sacrificing performance.
- 💰 **Cost Effective:** It's free! 🤑
- 🎯 **Diverse Use Cases:** Fully compatible with Diffusers 🧨 and the wider community.
- 🔧 **Developer-Friendly:** Simple requests, fast responses.
---
## Available Models
* **VAE Decode 🖼️:** Quickly decode latent representations into high-quality images without compromising performance or workflow speed.
* **VAE Encode 🔢:** Efficiently encode images into latent representations for generation and training.
* **Text Encoders 📃 (coming soon):** Compute text embeddings for your prompts quickly and accurately, ensuring a smooth and high-quality workflow.
---
## Integrations
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.
## Changelog
- March 10 2025: Added VAE encode
- March 2 2025: Initial release with VAE decoding
## Contents
The documentation is organized into three sections:
* **VAE Decode** Learn the basics of how to use VAE Decode with Hybrid Inference.
* **VAE Encode** Learn the basics of how to use VAE Encode with Hybrid Inference.
* **API Reference** Dive into task-specific settings and parameters.

View File

@@ -1,345 +0,0 @@
# Getting Started: VAE Decode with Hybrid Inference
VAE decode is an essential component of diffusion models - turning latent representations into images or videos.
## Memory
These tables demonstrate the VRAM requirements for VAE decode with SD v1 and SD XL on different GPUs.
For the majority of these GPUs the memory usage % dictates other models (text encoders, UNet/Transformer) must be offloaded, or tiled decoding has to be used which increases time taken and impacts quality.
<details><summary>SD v1.5</summary>
| GPU | Resolution | Time (seconds) | Memory (%) | Tiled Time (secs) | Tiled Memory (%) |
| --- | --- | --- | --- | --- | --- |
| NVIDIA GeForce RTX 4090 | 512x512 | 0.031 | 5.60% | 0.031 (0%) | 5.60% |
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.148 | 20.00% | 0.301 (+103%) | 5.60% |
| NVIDIA GeForce RTX 4080 | 512x512 | 0.05 | 8.40% | 0.050 (0%) | 8.40% |
| NVIDIA GeForce RTX 4080 | 1024x1024 | 0.224 | 30.00% | 0.356 (+59%) | 8.40% |
| NVIDIA GeForce RTX 4070 Ti | 512x512 | 0.066 | 11.30% | 0.066 (0%) | 11.30% |
| NVIDIA GeForce RTX 4070 Ti | 1024x1024 | 0.284 | 40.50% | 0.454 (+60%) | 11.40% |
| NVIDIA GeForce RTX 3090 | 512x512 | 0.062 | 5.20% | 0.062 (0%) | 5.20% |
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.253 | 18.50% | 0.464 (+83%) | 5.20% |
| NVIDIA GeForce RTX 3080 | 512x512 | 0.07 | 12.80% | 0.070 (0%) | 12.80% |
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.286 | 45.30% | 0.466 (+63%) | 12.90% |
| NVIDIA GeForce RTX 3070 | 512x512 | 0.102 | 15.90% | 0.102 (0%) | 15.90% |
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.421 | 56.30% | 0.746 (+77%) | 16.00% |
</details>
<details><summary>SDXL</summary>
| GPU | Resolution | Time (seconds) | Memory Consumed (%) | Tiled Time (seconds) | Tiled Memory (%) |
| --- | --- | --- | --- | --- | --- |
| NVIDIA GeForce RTX 4090 | 512x512 | 0.057 | 10.00% | 0.057 (0%) | 10.00% |
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.256 | 35.50% | 0.257 (+0.4%) | 35.50% |
| NVIDIA GeForce RTX 4080 | 512x512 | 0.092 | 15.00% | 0.092 (0%) | 15.00% |
| NVIDIA GeForce RTX 4080 | 1024x1024 | 0.406 | 53.30% | 0.406 (0%) | 53.30% |
| NVIDIA GeForce RTX 4070 Ti | 512x512 | 0.121 | 20.20% | 0.120 (-0.8%) | 20.20% |
| NVIDIA GeForce RTX 4070 Ti | 1024x1024 | 0.519 | 72.00% | 0.519 (0%) | 72.00% |
| NVIDIA GeForce RTX 3090 | 512x512 | 0.107 | 10.50% | 0.107 (0%) | 10.50% |
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.459 | 38.00% | 0.460 (+0.2%) | 38.00% |
| NVIDIA GeForce RTX 3080 | 512x512 | 0.121 | 25.60% | 0.121 (0%) | 25.60% |
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.524 | 93.00% | 0.524 (0%) | 93.00% |
| NVIDIA GeForce RTX 3070 | 512x512 | 0.183 | 31.80% | 0.183 (0%) | 31.80% |
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.794 | 96.40% | 0.794 (0%) | 96.40% |
</details>
## Available VAEs
| | **Endpoint** | **Model** |
|:-:|:-----------:|:--------:|
| **Stable Diffusion v1** | [https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud](https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud) | [`stabilityai/sd-vae-ft-mse`](https://hf.co/stabilityai/sd-vae-ft-mse) |
| **Stable Diffusion XL** | [https://x2dmsqunjd6k9prw.us-east-1.aws.endpoints.huggingface.cloud](https://x2dmsqunjd6k9prw.us-east-1.aws.endpoints.huggingface.cloud) | [`madebyollin/sdxl-vae-fp16-fix`](https://hf.co/madebyollin/sdxl-vae-fp16-fix) |
| **Flux** | [https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud](https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud) | [`black-forest-labs/FLUX.1-schnell`](https://hf.co/black-forest-labs/FLUX.1-schnell) |
| **HunyuanVideo** | [https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud](https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud) | [`hunyuanvideo-community/HunyuanVideo`](https://hf.co/hunyuanvideo-community/HunyuanVideo) |
> [!TIP]
> Model support can be requested [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
## Code
> [!TIP]
> Install `diffusers` from `main` to run the code: `pip install git+https://github.com/huggingface/diffusers@main`
A helper method simplifies interacting with Hybrid Inference.
```python
from diffusers.utils.remote_utils import remote_decode
```
### Basic example
Here, we show how to use the remote VAE on random tensors.
<details><summary>Code</summary>
```python
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=torch.randn([1, 4, 64, 64], dtype=torch.float16),
scaling_factor=0.18215,
)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/output.png"/>
</figure>
Usage for Flux is slightly different. Flux latents are packed so we need to send the `height` and `width`.
<details><summary>Code</summary>
```python
image = remote_decode(
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=torch.randn([1, 4096, 64], dtype=torch.float16),
height=1024,
width=1024,
scaling_factor=0.3611,
shift_factor=0.1159,
)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/flux_random_latent.png"/>
</figure>
Finally, an example for HunyuanVideo.
<details><summary>Code</summary>
```python
video = remote_decode(
endpoint="https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=torch.randn([1, 16, 3, 40, 64], dtype=torch.float16),
output_type="mp4",
)
with open("video.mp4", "wb") as f:
f.write(video)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<video
alt="queue.mp4"
autoplay loop autobuffer muted playsinline
>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/video_1.mp4" type="video/mp4">
</video>
</figure>
### Generation
But we want to use the VAE on an actual pipeline to get an actual image, not random noise. The example below shows how to do it with SD v1.5.
<details><summary>Code</summary>
```python
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
variant="fp16",
vae=None,
).to("cuda")
prompt = "Strawberry ice cream, in a stylish modern glass, coconut, splashing milk cream and honey, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious"
latent = pipe(
prompt=prompt,
output_type="latent",
).images
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
scaling_factor=0.18215,
)
image.save("test.jpg")
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/test.jpg"/>
</figure>
Heres another example with Flux.
<details><summary>Code</summary>
```python
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained(
"black-forest-labs/FLUX.1-schnell",
torch_dtype=torch.bfloat16,
vae=None,
).to("cuda")
prompt = "Strawberry ice cream, in a stylish modern glass, coconut, splashing milk cream and honey, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious"
latent = pipe(
prompt=prompt,
guidance_scale=0.0,
num_inference_steps=4,
output_type="latent",
).images
image = remote_decode(
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
height=1024,
width=1024,
scaling_factor=0.3611,
shift_factor=0.1159,
)
image.save("test.jpg")
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/test_1.jpg"/>
</figure>
Heres an example with HunyuanVideo.
<details><summary>Code</summary>
```python
from diffusers import HunyuanVideoPipeline, HunyuanVideoTransformer3DModel
model_id = "hunyuanvideo-community/HunyuanVideo"
transformer = HunyuanVideoTransformer3DModel.from_pretrained(
model_id, subfolder="transformer", torch_dtype=torch.bfloat16
)
pipe = HunyuanVideoPipeline.from_pretrained(
model_id, transformer=transformer, vae=None, torch_dtype=torch.float16
).to("cuda")
latent = pipe(
prompt="A cat walks on the grass, realistic",
height=320,
width=512,
num_frames=61,
num_inference_steps=30,
output_type="latent",
).frames
video = remote_decode(
endpoint="https://o7ywnmrahorts457.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
output_type="mp4",
)
if isinstance(video, bytes):
with open("video.mp4", "wb") as f:
f.write(video)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<video
alt="queue.mp4"
autoplay loop autobuffer muted playsinline
>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/video.mp4" type="video/mp4">
</video>
</figure>
### Queueing
One of the great benefits of using a remote VAE is that we can queue multiple generation requests. While the current latent is being processed for decoding, we can already queue another one. This helps improve concurrency.
<details><summary>Code</summary>
```python
import queue
import threading
from IPython.display import display
from diffusers import StableDiffusionPipeline
def decode_worker(q: queue.Queue):
while True:
item = q.get()
if item is None:
break
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=item,
scaling_factor=0.18215,
)
display(image)
q.task_done()
q = queue.Queue()
thread = threading.Thread(target=decode_worker, args=(q,), daemon=True)
thread.start()
def decode(latent: torch.Tensor):
q.put(latent)
prompts = [
"Blueberry ice cream, in a stylish modern glass , ice cubes, nuts, mint leaves, splashing milk cream, in a gradient purple background, fluid motion, dynamic movement, cinematic lighting, Mysterious",
"Lemonade in a glass, mint leaves, in an aqua and white background, flowers, ice cubes, halo, fluid motion, dynamic movement, soft lighting, digital painting, rule of thirds composition, Art by Greg rutkowski, Coby whitmore",
"Comic book art, beautiful, vintage, pastel neon colors, extremely detailed pupils, delicate features, light on face, slight smile, Artgerm, Mary Blair, Edmund Dulac, long dark locks, bangs, glowing, fashionable style, fairytale ambience, hot pink.",
"Masterpiece, vanilla cone ice cream garnished with chocolate syrup, crushed nuts, choco flakes, in a brown background, gold, cinematic lighting, Art by WLOP",
"A bowl of milk, falling cornflakes, berries, blueberries, in a white background, soft lighting, intricate details, rule of thirds, octane render, volumetric lighting",
"Cold Coffee with cream, crushed almonds, in a glass, choco flakes, ice cubes, wet, in a wooden background, cinematic lighting, hyper realistic painting, art by Carne Griffiths, octane render, volumetric lighting, fluid motion, dynamic movement, muted colors,",
]
pipe = StableDiffusionPipeline.from_pretrained(
"Lykon/dreamshaper-8",
torch_dtype=torch.float16,
vae=None,
).to("cuda")
pipe.unet = pipe.unet.to(memory_format=torch.channels_last)
pipe.unet = torch.compile(pipe.unet, mode="reduce-overhead", fullgraph=True)
_ = pipe(
prompt=prompts[0],
output_type="latent",
)
for prompt in prompts:
latent = pipe(
prompt=prompt,
output_type="latent",
).images
decode(latent)
q.put(None)
thread.join()
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<video
alt="queue.mp4"
autoplay loop autobuffer muted playsinline
>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/queue.mp4" type="video/mp4">
</video>
</figure>
## Integrations
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.

View File

@@ -1,183 +0,0 @@
# Getting Started: VAE Encode with Hybrid Inference
VAE encode is used for training, image-to-image and image-to-video - turning into images or videos into latent representations.
## Memory
These tables demonstrate the VRAM requirements for VAE encode with SD v1 and SD XL on different GPUs.
For the majority of these GPUs the memory usage % dictates other models (text encoders, UNet/Transformer) must be offloaded, or tiled encoding has to be used which increases time taken and impacts quality.
<details><summary>SD v1.5</summary>
| GPU | Resolution | Time (seconds) | Memory (%) | Tiled Time (secs) | Tiled Memory (%) |
|:------------------------------|:-------------|-----------------:|-------------:|--------------------:|-------------------:|
| NVIDIA GeForce RTX 4090 | 512x512 | 0.015 | 3.51901 | 0.015 | 3.51901 |
| NVIDIA GeForce RTX 4090 | 256x256 | 0.004 | 1.3154 | 0.005 | 1.3154 |
| NVIDIA GeForce RTX 4090 | 2048x2048 | 0.402 | 47.1852 | 0.496 | 3.51901 |
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.078 | 12.2658 | 0.094 | 3.51901 |
| NVIDIA GeForce RTX 4080 SUPER | 512x512 | 0.023 | 5.30105 | 0.023 | 5.30105 |
| NVIDIA GeForce RTX 4080 SUPER | 256x256 | 0.006 | 1.98152 | 0.006 | 1.98152 |
| NVIDIA GeForce RTX 4080 SUPER | 2048x2048 | 0.574 | 71.08 | 0.656 | 5.30105 |
| NVIDIA GeForce RTX 4080 SUPER | 1024x1024 | 0.111 | 18.4772 | 0.14 | 5.30105 |
| NVIDIA GeForce RTX 3090 | 512x512 | 0.032 | 3.52782 | 0.032 | 3.52782 |
| NVIDIA GeForce RTX 3090 | 256x256 | 0.01 | 1.31869 | 0.009 | 1.31869 |
| NVIDIA GeForce RTX 3090 | 2048x2048 | 0.742 | 47.3033 | 0.954 | 3.52782 |
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.136 | 12.2965 | 0.207 | 3.52782 |
| NVIDIA GeForce RTX 3080 | 512x512 | 0.036 | 8.51761 | 0.036 | 8.51761 |
| NVIDIA GeForce RTX 3080 | 256x256 | 0.01 | 3.18387 | 0.01 | 3.18387 |
| NVIDIA GeForce RTX 3080 | 2048x2048 | 0.863 | 86.7424 | 1.191 | 8.51761 |
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.157 | 29.6888 | 0.227 | 8.51761 |
| NVIDIA GeForce RTX 3070 | 512x512 | 0.051 | 10.6941 | 0.051 | 10.6941 |
| NVIDIA GeForce RTX 3070 | 256x256 | 0.015 | 3.99743 | 0.015 | 3.99743 |
| NVIDIA GeForce RTX 3070 | 2048x2048 | 1.217 | 96.054 | 1.482 | 10.6941 |
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.223 | 37.2751 | 0.327 | 10.6941 |
</details>
<details><summary>SDXL</summary>
| GPU | Resolution | Time (seconds) | Memory Consumed (%) | Tiled Time (seconds) | Tiled Memory (%) |
|:------------------------------|:-------------|-----------------:|----------------------:|-----------------------:|-------------------:|
| NVIDIA GeForce RTX 4090 | 512x512 | 0.029 | 4.95707 | 0.029 | 4.95707 |
| NVIDIA GeForce RTX 4090 | 256x256 | 0.007 | 2.29666 | 0.007 | 2.29666 |
| NVIDIA GeForce RTX 4090 | 2048x2048 | 0.873 | 66.3452 | 0.863 | 15.5649 |
| NVIDIA GeForce RTX 4090 | 1024x1024 | 0.142 | 15.5479 | 0.143 | 15.5479 |
| NVIDIA GeForce RTX 4080 SUPER | 512x512 | 0.044 | 7.46735 | 0.044 | 7.46735 |
| NVIDIA GeForce RTX 4080 SUPER | 256x256 | 0.01 | 3.4597 | 0.01 | 3.4597 |
| NVIDIA GeForce RTX 4080 SUPER | 2048x2048 | 1.317 | 87.1615 | 1.291 | 23.447 |
| NVIDIA GeForce RTX 4080 SUPER | 1024x1024 | 0.213 | 23.4215 | 0.214 | 23.4215 |
| NVIDIA GeForce RTX 3090 | 512x512 | 0.058 | 5.65638 | 0.058 | 5.65638 |
| NVIDIA GeForce RTX 3090 | 256x256 | 0.016 | 2.45081 | 0.016 | 2.45081 |
| NVIDIA GeForce RTX 3090 | 2048x2048 | 1.755 | 77.8239 | 1.614 | 18.4193 |
| NVIDIA GeForce RTX 3090 | 1024x1024 | 0.265 | 18.4023 | 0.265 | 18.4023 |
| NVIDIA GeForce RTX 3080 | 512x512 | 0.064 | 13.6568 | 0.064 | 13.6568 |
| NVIDIA GeForce RTX 3080 | 256x256 | 0.018 | 5.91728 | 0.018 | 5.91728 |
| NVIDIA GeForce RTX 3080 | 2048x2048 | OOM | OOM | 1.866 | 44.4717 |
| NVIDIA GeForce RTX 3080 | 1024x1024 | 0.302 | 44.4308 | 0.302 | 44.4308 |
| NVIDIA GeForce RTX 3070 | 512x512 | 0.093 | 17.1465 | 0.093 | 17.1465 |
| NVIDIA GeForce RTX 3070 | 256x256 | 0.025 | 7.42931 | 0.026 | 7.42931 |
| NVIDIA GeForce RTX 3070 | 2048x2048 | OOM | OOM | 2.674 | 55.8355 |
| NVIDIA GeForce RTX 3070 | 1024x1024 | 0.443 | 55.7841 | 0.443 | 55.7841 |
</details>
## Available VAEs
| | **Endpoint** | **Model** |
|:-:|:-----------:|:--------:|
| **Stable Diffusion v1** | [https://qc6479g0aac6qwy9.us-east-1.aws.endpoints.huggingface.cloud](https://qc6479g0aac6qwy9.us-east-1.aws.endpoints.huggingface.cloud) | [`stabilityai/sd-vae-ft-mse`](https://hf.co/stabilityai/sd-vae-ft-mse) |
| **Stable Diffusion XL** | [https://xjqqhmyn62rog84g.us-east-1.aws.endpoints.huggingface.cloud](https://xjqqhmyn62rog84g.us-east-1.aws.endpoints.huggingface.cloud) | [`madebyollin/sdxl-vae-fp16-fix`](https://hf.co/madebyollin/sdxl-vae-fp16-fix) |
| **Flux** | [https://ptccx55jz97f9zgo.us-east-1.aws.endpoints.huggingface.cloud](https://ptccx55jz97f9zgo.us-east-1.aws.endpoints.huggingface.cloud) | [`black-forest-labs/FLUX.1-schnell`](https://hf.co/black-forest-labs/FLUX.1-schnell) |
> [!TIP]
> Model support can be requested [here](https://github.com/huggingface/diffusers/issues/new?template=remote-vae-pilot-feedback.yml).
## Code
> [!TIP]
> Install `diffusers` from `main` to run the code: `pip install git+https://github.com/huggingface/diffusers@main`
A helper method simplifies interacting with Hybrid Inference.
```python
from diffusers.utils.remote_utils import remote_encode
```
### Basic example
Let's encode an image, then decode it to demonstrate.
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg"/>
</figure>
<details><summary>Code</summary>
```python
from diffusers.utils import load_image
from diffusers.utils.remote_utils import remote_decode
image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/astronaut.jpg?download=true")
latent = remote_encode(
endpoint="https://ptccx55jz97f9zgo.us-east-1.aws.endpoints.huggingface.cloud/",
scaling_factor=0.3611,
shift_factor=0.1159,
)
decoded = remote_decode(
endpoint="https://whhx50ex1aryqvw6.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
scaling_factor=0.3611,
shift_factor=0.1159,
)
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/decoded.png"/>
</figure>
### Generation
Now let's look at a generation example, we'll encode the image, generate then remotely decode too!
<details><summary>Code</summary>
```python
import torch
from diffusers import StableDiffusionImg2ImgPipeline
from diffusers.utils import load_image
from diffusers.utils.remote_utils import remote_decode, remote_encode
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
variant="fp16",
vae=None,
).to("cuda")
init_image = load_image(
"https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
)
init_image = init_image.resize((768, 512))
init_latent = remote_encode(
endpoint="https://qc6479g0aac6qwy9.us-east-1.aws.endpoints.huggingface.cloud/",
image=init_image,
scaling_factor=0.18215,
)
prompt = "A fantasy landscape, trending on artstation"
latent = pipe(
prompt=prompt,
image=init_latent,
strength=0.75,
output_type="latent",
).images
image = remote_decode(
endpoint="https://q1bj3bpq6kzilnsu.us-east-1.aws.endpoints.huggingface.cloud/",
tensor=latent,
scaling_factor=0.18215,
)
image.save("fantasy_landscape.jpg")
```
</details>
<figure class="image flex flex-col items-center justify-center text-center m-0 w-full">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/blog/remote_vae/fantasy_landscape.png"/>
</figure>
## Integrations
* **[SD.Next](https://github.com/vladmandic/sdnext):** All-in-one UI with direct supports Hybrid Inference.
* **[ComfyUI-HFRemoteVae](https://github.com/kijai/ComfyUI-HFRemoteVae):** ComfyUI node for Hybrid Inference.

View File

@@ -23,60 +23,32 @@ You should install 🤗 Diffusers in a [virtual environment](https://docs.python
If you're unfamiliar with Python virtual environments, take a look at this [guide](https://packaging.python.org/guides/installing-using-pip-and-virtual-environments/).
A virtual environment makes it easier to manage different projects and avoid compatibility issues between dependencies.
Create a virtual environment with Python or [uv](https://docs.astral.sh/uv/) (refer to [Installation](https://docs.astral.sh/uv/getting-started/installation/) for installation instructions), a fast Rust-based Python package and project manager.
<hfoptions id="install">
<hfoption id="uv">
Start by creating a virtual environment in your project directory:
```bash
uv venv my-env
source my-env/bin/activate
python -m venv .env
```
</hfoption>
<hfoption id="Python">
Activate the virtual environment:
```bash
python -m venv my-env
source my-env/bin/activate
source .env/bin/activate
```
</hfoption>
</hfoptions>
You should also install 🤗 Transformers because 🤗 Diffusers relies on its models.
You should also install 🤗 Transformers because 🤗 Diffusers relies on its models:
<frameworkcontent>
<pt>
PyTorch only supports Python 3.8 - 3.11 on Windows. Install Diffusers with uv.
```bash
uv install diffusers["torch"] transformers
```
You can also install Diffusers with pip.
Note - PyTorch only supports Python 3.8 - 3.11 on Windows.
```bash
pip install diffusers["torch"] transformers
```
</pt>
<jax>
Install Diffusers with uv.
```bash
uv pip install diffusers["flax"] transformers
```
You can also install Diffusers with pip.
```bash
pip install diffusers["flax"] transformers
```
</jax>
</frameworkcontent>
@@ -161,10 +133,10 @@ Your Python environment will find the `main` version of 🤗 Diffusers on the ne
Model weights and files are downloaded from the Hub to a cache which is usually your home directory. You can change the cache location by specifying the `HF_HOME` or `HUGGINFACE_HUB_CACHE` environment variables or configuring the `cache_dir` parameter in methods like [`~DiffusionPipeline.from_pretrained`].
Cached files allow you to run 🤗 Diffusers offline. To prevent 🤗 Diffusers from connecting to the internet, set the `HF_HUB_OFFLINE` environment variable to `1` and 🤗 Diffusers will only load previously downloaded files in the cache.
Cached files allow you to run 🤗 Diffusers offline. To prevent 🤗 Diffusers from connecting to the internet, set the `HF_HUB_OFFLINE` environment variable to `True` and 🤗 Diffusers will only load previously downloaded files in the cache.
```shell
export HF_HUB_OFFLINE=1
export HF_HUB_OFFLINE=True
```
For more details about managing and cleaning the cache, take a look at the [caching](https://huggingface.co/docs/huggingface_hub/guides/manage-cache) guide.
@@ -179,16 +151,14 @@ Telemetry is only sent when loading models and pipelines from the Hub,
and it is not collected if you're loading local files.
We understand that not everyone wants to share additional information,and we respect your privacy.
You can disable telemetry collection by setting the `HF_HUB_DISABLE_TELEMETRY` environment variable from your terminal:
You can disable telemetry collection by setting the `DISABLE_TELEMETRY` environment variable from your terminal:
On Linux/MacOS:
```bash
export HF_HUB_DISABLE_TELEMETRY=1
export DISABLE_TELEMETRY=YES
```
On Windows:
```bash
set HF_HUB_DISABLE_TELEMETRY=1
set DISABLE_TELEMETRY=YES
```

View File

@@ -10,211 +10,120 @@ an "AS IS" BASIS, WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express o
specific language governing permissions and limitations under the License.
-->
# Accelerate inference
# Speed up inference
Diffusion models are slow at inference because generation is an iterative process where noise is gradually refined into an image or video over a certain number of "steps". To speedup this process, you can try experimenting with different [schedulers](../api/schedulers/overview), reduce the precision of the model weights for faster computations, use more memory-efficient attention mechanisms, and more.
There are several ways to optimize Diffusers for inference speed, such as reducing the computational burden by lowering the data precision or using a lightweight distilled model. There are also memory-efficient attention implementations, [xFormers](xformers) and [scaled dot product attention](https://pytorch.org/docs/stable/generated/torch.nn.functional.scaled_dot_product_attention.html) in PyTorch 2.0, that reduce memory usage which also indirectly speeds up inference. Different speed optimizations can be stacked together to get the fastest inference times.
Combine and use these techniques together to make inference faster than using any single technique on its own.
> [!TIP]
> Optimizing for inference speed or reduced memory usage can lead to improved performance in the other category, so you should try to optimize for both whenever you can. This guide focuses on inference speed, but you can learn more about lowering memory usage in the [Reduce memory usage](memory) guide.
This guide will go over how to accelerate inference.
The inference times below are obtained from generating a single 512x512 image from the prompt "a photo of an astronaut riding a horse on mars" with 50 DDIM steps on a NVIDIA A100.
## Model data type
| setup | latency | speed-up |
|----------|---------|----------|
| baseline | 5.27s | x1 |
| tf32 | 4.14s | x1.27 |
| fp16 | 3.51s | x1.50 |
| combined | 3.41s | x1.54 |
The precision and data type of the model weights affect inference speed because a higher precision requires more memory to load and more time to perform the computations. PyTorch loads model weights in float32 or full precision by default, so changing the data type is a simple way to quickly get faster inference.
## TensorFloat-32
<hfoptions id="dtypes">
<hfoption id="bfloat16">
On Ampere and later CUDA devices, matrix multiplications and convolutions can use the [TensorFloat-32 (tf32)](https://blogs.nvidia.com/blog/2020/05/14/tensorfloat-32-precision-format/) mode for faster, but slightly less accurate computations. By default, PyTorch enables tf32 mode for convolutions but not matrix multiplications. Unless your network requires full float32 precision, we recommend enabling tf32 for matrix multiplications. It can significantly speed up computations with typically negligible loss in numerical accuracy.
bfloat16 is similar to float16 but it is more robust to numerical errors. Hardware support for bfloat16 varies, but most modern GPUs are capable of supporting bfloat16.
```py
```python
import torch
from diffusers import StableDiffusionXLPipeline
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.bfloat16
).to("cuda")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
pipeline(prompt, num_inference_steps=30).images[0]
```
</hfoption>
<hfoption id="float16">
float16 is similar to bfloat16 but may be more prone to numerical errors.
```py
import torch
from diffusers import StableDiffusionXLPipeline
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16
).to("cuda")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
pipeline(prompt, num_inference_steps=30).images[0]
```
</hfoption>
<hfoption id="TensorFloat-32">
[TensorFloat-32 (tf32)](https://blogs.nvidia.com/blog/2020/05/14/tensorfloat-32-precision-format/) mode is supported on NVIDIA Ampere GPUs and it computes the convolution and matrix multiplication operations in tf32. Storage and other operations are kept in float32. This enables significantly faster computations when combined with bfloat16 or float16.
PyTorch only enables tf32 mode for convolutions by default and you'll need to explicitly enable it for matrix multiplications.
```py
import torch
from diffusers import StableDiffusionXLPipeline
torch.backends.cuda.matmul.allow_tf32 = True
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.bfloat16
).to("cuda")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
pipeline(prompt, num_inference_steps=30).images[0]
```
Refer to the [mixed precision training](https://huggingface.co/docs/transformers/en/perf_train_gpu_one#mixed-precision) docs for more details.
Learn more about tf32 in the [Mixed precision training](https://huggingface.co/docs/transformers/en/perf_train_gpu_one#tf32) guide.
</hfoption>
</hfoptions>
## Half-precision weights
## Scaled dot product attention
To save GPU memory and get more speed, set `torch_dtype=torch.float16` to load and run the model weights directly with half-precision weights.
> [!TIP]
> Memory-efficient attention optimizes for inference speed *and* [memory usage](./memory#memory-efficient-attention)!
[Scaled dot product attention (SDPA)](https://pytorch.org/docs/stable/generated/torch.nn.functional.scaled_dot_product_attention.html) implements several attention backends, [FlashAttention](https://github.com/Dao-AILab/flash-attention), [xFormers](https://github.com/facebookresearch/xformers), and a native C++ implementation. It automatically selects the most optimal backend for your hardware.
SDPA is enabled by default if you're using PyTorch >= 2.0 and no additional changes are required to your code. You could try experimenting with other attention backends though if you'd like to choose your own. The example below uses the [torch.nn.attention.sdpa_kernel](https://pytorch.org/docs/stable/generated/torch.nn.attention.sdpa_kernel.html) context manager to enable efficient attention.
```py
from torch.nn.attention import SDPBackend, sdpa_kernel
```Python
import torch
from diffusers import StableDiffusionXLPipeline
from diffusers import DiffusionPipeline
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.bfloat16
).to("cuda")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
with sdpa_kernel(SDPBackend.EFFICIENT_ATTENTION):
image = pipeline(prompt, num_inference_steps=30).images[0]
```
## torch.compile
[torch.compile](https://pytorch.org/tutorials/intermediate/torch_compile_tutorial.html) accelerates inference by compiling PyTorch code and operations into optimized kernels. Diffusers typically compiles the more compute-intensive models like the UNet, transformer, or VAE.
Enable the following compiler settings for maximum speed (refer to the [full list](https://github.com/pytorch/pytorch/blob/main/torch/_inductor/config.py) for more options).
```py
import torch
from diffusers import StableDiffusionXLPipeline
torch._inductor.config.conv_1x1_as_mm = True
torch._inductor.config.coordinate_descent_tuning = True
torch._inductor.config.epilogue_fusion = False
torch._inductor.config.coordinate_descent_check_all_directions = True
```
Load and compile the UNet and VAE. There are several different modes you can choose from, but `"max-autotune"` optimizes for the fastest speed by compiling to a CUDA graph. CUDA graphs effectively reduces the overhead by launching multiple GPU operations through a single CPU operation.
> [!TIP]
> With PyTorch 2.3.1, you can control the caching behavior of torch.compile. This is particularly beneficial for compilation modes like `"max-autotune"` which performs a grid-search over several compilation flags to find the optimal configuration. Learn more in the [Compile Time Caching in torch.compile](https://pytorch.org/tutorials/recipes/torch_compile_caching_tutorial.html) tutorial.
Changing the memory layout to [channels_last](./memory#torchchannels_last) also optimizes memory and inference speed.
```py
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16
).to("cuda")
pipeline.unet.to(memory_format=torch.channels_last)
pipeline.vae.to(memory_format=torch.channels_last)
pipeline.unet = torch.compile(
pipeline.unet, mode="max-autotune", fullgraph=True
pipe = DiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
)
pipeline.vae.decode = torch.compile(
pipeline.vae.decode,
mode="max-autotune",
fullgraph=True
)
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
pipeline(prompt, num_inference_steps=30).images[0]
pipe = pipe.to("cuda")
```
Compilation is slow the first time, but once compiled, it is significantly faster. Try to only use the compiled pipeline on the same type of inference operations. Calling the compiled pipeline on a different image size retriggers compilation which is slow and inefficient.
### Graph breaks
It is important to specify `fullgraph=True` in torch.compile to ensure there are no graph breaks in the underlying model. This allows you to take advantage of torch.compile without any performance degradation. For the UNet and VAE, this changes how you access the return variables.
```diff
- latents = unet(
- latents, timestep=timestep, encoder_hidden_states=prompt_embeds
-).sample
+ latents = unet(
+ latents, timestep=timestep, encoder_hidden_states=prompt_embeds, return_dict=False
+)[0]
```
### GPU sync
The `step()` function is [called](https://github.com/huggingface/diffusers/blob/1d686bac8146037e97f3fd8c56e4063230f71751/src/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl.py#L1228) on the scheduler each time after the denoiser makes a prediction, and the `sigmas` variable is [indexed](https://github.com/huggingface/diffusers/blob/1d686bac8146037e97f3fd8c56e4063230f71751/src/diffusers/schedulers/scheduling_euler_discrete.py#L476). When placed on the GPU, it introduces latency because of the communication sync between the CPU and GPU. It becomes more evident when the denoiser has already been compiled.
In general, the `sigmas` should [stay on the CPU](https://github.com/huggingface/diffusers/blob/35a969d297cba69110d175ee79c59312b9f49e1e/src/diffusers/schedulers/scheduling_euler_discrete.py#L240) to avoid the communication sync and latency.
## Dynamic quantization
[Dynamic quantization](https://pytorch.org/tutorials/recipes/recipes/dynamic_quantization.html) improves inference speed by reducing precision to enable faster math operations. This particular type of quantization determines how to scale the activations based on the data at runtime rather than using a fixed scaling factor. As a result, the scaling factor is more accurately aligned with the data.
The example below applies [dynamic int8 quantization](https://pytorch.org/tutorials/recipes/recipes/dynamic_quantization.html) to the UNet and VAE with the [torchao](../quantization/torchao) library.
> [!TIP]
> Refer to our [torchao](../quantization/torchao) docs to learn more about how to use the Diffusers torchao integration.
Configure the compiler tags for maximum speed.
```py
import torch
from torchao import apply_dynamic_quant
from diffusers import StableDiffusionXLPipeline
torch._inductor.config.conv_1x1_as_mm = True
torch._inductor.config.coordinate_descent_tuning = True
torch._inductor.config.epilogue_fusion = False
torch._inductor.config.coordinate_descent_check_all_directions = True
torch._inductor.config.force_fuse_int_mm_with_mul = True
torch._inductor.config.use_mixed_mm = True
```
Filter out some linear layers in the UNet and VAE which don't benefit from dynamic quantization with the [dynamic_quant_filter_fn](https://github.com/huggingface/diffusion-fast/blob/0f169640b1db106fe6a479f78c1ed3bfaeba3386/utils/pipeline_utils.py#L16).
```py
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.bfloat16
).to("cuda")
apply_dynamic_quant(pipeline.unet, dynamic_quant_filter_fn)
apply_dynamic_quant(pipeline.vae, dynamic_quant_filter_fn)
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
pipeline(prompt, num_inference_steps=30).images[0]
```
## Fused projection matrices
> [!WARNING]
> The [fuse_qkv_projections](https://github.com/huggingface/diffusers/blob/58431f102cf39c3c8a569f32d71b2ea8caa461e1/src/diffusers/pipelines/pipeline_utils.py#L2034) method is experimental and support is limited to mostly Stable Diffusion pipelines. Take a look at this [PR](https://github.com/huggingface/diffusers/pull/6179) to learn more about how to enable it for other pipelines
> Don't use [torch.autocast](https://pytorch.org/docs/stable/amp.html#torch.autocast) in any of the pipelines as it can lead to black images and is always slower than pure float16 precision.
An input is projected into three subspaces, represented by the projection matrices Q, K, and V, in an attention block. These projections are typically calculated separately, but you can horizontally combine these into a single matrix and perform the projection in a single step. It increases the size of the matrix multiplications of the input projections and also improves the impact of quantization.
## Distilled model
You could also use a distilled Stable Diffusion model and autoencoder to speed up inference. During distillation, many of the UNet's residual and attention blocks are shed to reduce the model size by 51% and improve latency on CPU/GPU by 43%. The distilled model is faster and uses less memory while generating images of comparable quality to the full Stable Diffusion model.
> [!TIP]
> Read the [Open-sourcing Knowledge Distillation Code and Weights of SD-Small and SD-Tiny](https://huggingface.co/blog/sd_distillation) blog post to learn more about how knowledge distillation training works to produce a faster, smaller, and cheaper generative model.
The inference times below are obtained from generating 4 images from the prompt "a photo of an astronaut riding a horse on mars" with 25 PNDM steps on a NVIDIA A100. Each generation is repeated 3 times with the distilled Stable Diffusion v1.4 model by [Nota AI](https://hf.co/nota-ai).
| setup | latency | speed-up |
|------------------------------|---------|----------|
| baseline | 6.37s | x1 |
| distilled | 4.18s | x1.52 |
| distilled + tiny autoencoder | 3.83s | x1.66 |
Let's load the distilled Stable Diffusion model and compare it against the original Stable Diffusion model.
```py
pipeline.fuse_qkv_projections()
```
from diffusers import StableDiffusionPipeline
import torch
distilled = StableDiffusionPipeline.from_pretrained(
"nota-ai/bk-sdm-small", torch_dtype=torch.float16, use_safetensors=True,
).to("cuda")
prompt = "a golden vase with different flowers"
generator = torch.manual_seed(2023)
image = distilled("a golden vase with different flowers", num_inference_steps=25, generator=generator).images[0]
image
```
<div class="flex gap-4">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/original_sd.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">original Stable Diffusion</figcaption>
</div>
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/distilled_sd.png"/>
<figcaption class="mt-2 text-center text-sm text-gray-500">distilled Stable Diffusion</figcaption>
</div>
</div>
### Tiny AutoEncoder
To speed inference up even more, replace the autoencoder with a [distilled version](https://huggingface.co/sayakpaul/taesdxl-diffusers) of it.
```py
import torch
from diffusers import AutoencoderTiny, StableDiffusionPipeline
distilled = StableDiffusionPipeline.from_pretrained(
"nota-ai/bk-sdm-small", torch_dtype=torch.float16, use_safetensors=True,
).to("cuda")
distilled.vae = AutoencoderTiny.from_pretrained(
"sayakpaul/taesd-diffusers", torch_dtype=torch.float16, use_safetensors=True,
).to("cuda")
prompt = "a golden vase with different flowers"
generator = torch.manual_seed(2023)
image = distilled("a golden vase with different flowers", num_inference_steps=25, generator=generator).images[0]
image
```
<div class="flex justify-center">
<div>
<img class="rounded-xl" src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/distilled_sd_vae.png" />
<figcaption class="mt-2 text-center text-sm text-gray-500">distilled Stable Diffusion + Tiny AutoEncoder</figcaption>
</div>
</div>
More tiny autoencoder models for other Stable Diffusion models, like Stable Diffusion 3, are available from [madebyollin](https://huggingface.co/madebyollin).

View File

@@ -12,371 +12,173 @@ specific language governing permissions and limitations under the License.
# Reduce memory usage
Modern diffusion models like [Flux](../api/pipelines/flux) and [Wan](../api/pipelines/wan) have billions of parameters that take up a lot of memory on your hardware for inference. This is challenging because common GPUs often don't have sufficient memory. To overcome the memory limitations, you can use more than one GPU (if available), offload some of the pipeline components to the CPU, and more.
A barrier to using diffusion models is the large amount of memory required. To overcome this challenge, there are several memory-reducing techniques you can use to run even some of the largest models on free-tier or consumer GPUs. Some of these techniques can even be combined to further reduce memory usage.
This guide will show you how to reduce your memory usage.
<Tip>
> [!TIP]
> Keep in mind these techniques may need to be adjusted depending on the model! For example, a transformer-based diffusion model may not benefit equally from these inference speed optimizations as a UNet-based model.
In many cases, optimizing for memory or speed leads to improved performance in the other, so you should try to optimize for both whenever you can. This guide focuses on minimizing memory usage, but you can also learn more about how to [Speed up inference](fp16).
## Multiple GPUs
</Tip>
If you have access to more than one GPU, there a few options for efficiently loading and distributing a large model across your hardware. These features are supported by the [Accelerate](https://huggingface.co/docs/accelerate/index) library, so make sure it is installed first.
The results below are obtained from generating a single 512x512 image from the prompt a photo of an astronaut riding a horse on mars with 50 DDIM steps on a Nvidia Titan RTX, demonstrating the speed-up you can expect as a result of reduced memory consumption.
```bash
pip install -U accelerate
```
| | latency | speed-up |
| ---------------- | ------- | ------- |
| original | 9.50s | x1 |
| fp16 | 3.61s | x2.63 |
| channels last | 3.30s | x2.88 |
| traced UNet | 3.21s | x2.96 |
| memory-efficient attention | 2.63s | x3.61 |
### Sharded checkpoints
## Sliced VAE
Loading large checkpoints in several shards in useful because the shards are loaded one at a time. This keeps memory usage low, only requiring enough memory for the model size and the largest shard size. We recommend sharding when the fp32 checkpoint is greater than 5GB. The default shard size is 5GB.
Sliced VAE enables decoding large batches of images with limited VRAM or batches with 32 images or more by decoding the batches of latents one image at a time. You'll likely want to couple this with [`~ModelMixin.enable_xformers_memory_efficient_attention`] to reduce memory use further if you have xFormers installed.
Shard a checkpoint in [`~DiffusionPipeline.save_pretrained`] with the `max_shard_size` parameter.
```py
from diffusers import AutoModel
unet = AutoModel.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", subfolder="unet"
)
unet.save_pretrained("sdxl-unet-sharded", max_shard_size="5GB")
```
Now you can use the sharded checkpoint, instead of the regular checkpoint, to save memory.
```py
import torch
from diffusers import AutoModel, StableDiffusionXLPipeline
unet = AutoModel.from_pretrained(
"username/sdxl-unet-sharded", torch_dtype=torch.float16
)
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
unet=unet,
torch_dtype=torch.float16
).to("cuda")
```
### Device placement
> [!WARNING]
> Device placement is an experimental feature and the API may change. Only the `balanced` strategy is supported at the moment. We plan to support additional mapping strategies in the future.
The `device_map` parameter controls how the model components in a pipeline are distributed across devices. The `balanced` device placement strategy evenly splits the pipeline across all available devices.
```py
import torch
from diffusers import AutoModel, StableDiffusionXLPipeline
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
device_map="balanced"
)
```
You can inspect a pipeline's device map with `hf_device_map`.
```py
print(pipeline.hf_device_map)
{'unet': 1, 'vae': 1, 'safety_checker': 0, 'text_encoder': 0}
```
The `device_map` parameter also works on the model-level. This is useful for loading large models, such as the Flux diffusion transformer which has 12.5B parameters. Instead of `balanced`, set it to `"auto"` to automatically distribute a model across the fastest device first before moving to slower devices. Refer to the [Model sharding](../training/distributed_inference#model-sharding) docs for more details.
```py
import torch
from diffusers import AutoModel
transformer = AutoModel.from_pretrained(
"black-forest-labs/FLUX.1-dev",
subfolder="transformer",
device_map="auto",
torch_dtype=torch.bfloat16
)
```
For more fine-grained control, pass a dictionary to enforce the maximum GPU memory to use on each device. If a device is not in `max_memory`, it is ignored and pipeline components won't be distributed to it.
```py
import torch
from diffusers import AutoModel, StableDiffusionXLPipeline
max_memory = {0:"1GB", 1:"1GB"}
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
device_map="balanced",
max_memory=max_memory
)
```
Diffusers uses the maxmium memory of all devices by default, but if they don't fit on the GPUs, then you'll need to use a single GPU and offload to the CPU with the methods below.
- [`~DiffusionPipeline.enable_model_cpu_offload`] only works on a single GPU but a very large model may not fit on it
- [`~DiffusionPipeline.enable_sequential_cpu_offload`] may work but it is extremely slow and also limited to a single GPU
Use the [`~DiffusionPipeline.reset_device_map`] method to reset the `device_map`. This is necessary if you want to use methods like `.to()`, [`~DiffusionPipeline.enable_sequential_cpu_offload`], and [`~DiffusionPipeline.enable_model_cpu_offload`] on a pipeline that was device-mapped.
```py
pipeline.reset_device_map()
```
## VAE slicing
VAE slicing saves memory by splitting large batches of inputs into a single batch of data and separately processing them. This method works best when generating more than one image at a time.
For example, if you're generating 4 images at once, decoding would increase peak activation memory by 4x. VAE slicing reduces this by only decoding 1 image at a time instead of all 4 images at once.
Call [`~StableDiffusionPipeline.enable_vae_slicing`] to enable sliced VAE. You can expect a small increase in performance when decoding multi-image batches and no performance impact for single-image batches.
```py
import torch
from diffusers import AutoModel, StableDiffusionXLPipeline
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
torch_dtype=torch.float16,
).to("cuda")
pipeline.enable_vae_slicing()
pipeline(["An astronaut riding a horse on Mars"]*32).images[0]
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
```
> [!WARNING]
> [`AutoencoderKLWan`] and [`AsymmetricAutoencoderKL`] don't support slicing.
## VAE tiling
VAE tiling saves memory by dividing an image into smaller overlapping tiles instead of processing the entire image at once. This also reduces peak memory usage because the GPU is only processing a tile at a time.
Call [`~StableDiffusionPipeline.enable_vae_tiling`] to enable VAE tiling. The generated image may have some tone variation from tile-to-tile because they're decoded separately, but there shouldn't be any obvious seams between the tiles. Tiling is disabled for resolutions lower than a pre-specified (but configurable) limit. For example, this limit is 512x512 for the VAE in [`StableDiffusionPipeline`].
```py
import torch
from diffusers import AutoPipelineForImage2Image
from diffusers.utils import load_image
pipeline = AutoPipelineForImage2Image.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0", torch_dtype=torch.float16
).to("cuda")
pipeline.enable_vae_tiling()
init_image = load_image("https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/img2img-sdxl-init.png")
prompt = "Astronaut in a jungle, cold color palette, muted colors, detailed, 8k"
pipeline(prompt, image=init_image, strength=0.5).images[0]
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
```
> [!WARNING]
> [`AutoencoderKLWan`] and [`AsymmetricAutoencoderKL`] don't support tiling.
## CPU offloading
CPU offloading selectively moves weights from the GPU to the CPU. When a component is required, it is transferred to the GPU and when it isn't required, it is moved to the CPU. This method works on submodules rather than whole models. It saves memory by avoiding storing the entire model on the GPU.
CPU offloading dramatically reduces memory usage, but it is also **extremely slow** because submodules are passed back and forth multiple times between devices. It can often be impractical due to how slow it is.
> [!WARNING]
> Don't move the pipeline to CUDA before calling [`~DiffusionPipeline.enable_sequential_cpu_offload`], otherwise the amount of memory saved is only minimal (refer to this [issue](https://github.com/huggingface/diffusers/issues/1934) for more details). This is a stateful operation that installs hooks on the model.
Call [`~DiffusionPipeline.enable_sequential_cpu_offload`] to enable it on a pipeline.
```py
import torch
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"black-forest-labs/FLUX.1-schnell", torch_dtype=torch.bfloat16
)
pipeline.enable_sequential_cpu_offload()
pipeline(
prompt="An astronaut riding a horse on Mars",
guidance_scale=0.,
height=768,
width=1360,
num_inference_steps=4,
max_sequence_length=256,
).images[0]
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
```
## Model offloading
Model offloading moves entire models to the GPU instead of selectively moving *some* layers or model components. One of the main pipeline models, usually the text encoder, UNet, and VAE, is placed on the GPU while the other components are held on the CPU. Components like the UNet that run multiple times stays on the GPU until its completely finished and no longer needed. This eliminates the communication overhead of [CPU offloading](#cpu-offloading) and makes model offloading a faster alternative. The tradeoff is memory savings won't be as large.
> [!WARNING]
> Keep in mind that if models are reused outside the pipeline after hookes have been installed (see [Removing Hooks](https://huggingface.co/docs/accelerate/en/package_reference/big_modeling#accelerate.hooks.remove_hook_from_module) for more details), you need to run the entire pipeline and models in the expected order to properly offload them. This is a stateful operation that installs hooks on the model.
Call [`~DiffusionPipeline.enable_model_cpu_offload`] to enable it on a pipeline.
```py
import torch
from diffusers import DiffusionPipeline
pipeline = DiffusionPipeline.from_pretrained(
"black-forest-labs/FLUX.1-schnell", torch_dtype=torch.bfloat16
)
pipline.enable_model_cpu_offload()
pipeline(
prompt="An astronaut riding a horse on Mars",
guidance_scale=0.,
height=768,
width=1360,
num_inference_steps=4,
max_sequence_length=256,
).images[0]
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
```
[`~DiffusionPipeline.enable_model_cpu_offload`] also helps when you're using the [`~StableDiffusionXLPipeline.encode_prompt`] method on its own to generate the text encoders hidden state.
## Group offloading
Group offloading moves groups of internal layers ([torch.nn.ModuleList](https://pytorch.org/docs/stable/generated/torch.nn.ModuleList.html) or [torch.nn.Sequential](https://pytorch.org/docs/stable/generated/torch.nn.Sequential.html)) to the CPU. It uses less memory than [model offloading](#model-offloading) and it is faster than [CPU offloading](#cpu-offloading) because it reduces communication overhead.
> [!WARNING]
> Group offloading may not work with all models if the forward implementation contains weight-dependent device casting of inputs because it may clash with group offloading's device casting mechanism.
Call [`~ModelMixin.enable_group_offload`] to enable it for standard Diffusers model components that inherit from [`ModelMixin`]. For other model components that don't inherit from [`ModelMixin`], such as a generic [torch.nn.Module](https://pytorch.org/docs/stable/generated/torch.nn.Module.html), use [`~hooks.apply_group_offloading`] instead.
The `offload_type` parameter can be set to `block_level` or `leaf_level`.
- `block_level` offloads groups of layers based on the `num_blocks_per_group` parameter. For example, if `num_blocks_per_group=2` on a model with 40 layers, 2 layers are onloaded and offloaded at a time (20 total onloads/offloads). This drastically reduces memory requirements.
- `leaf_level` offloads individual layers at the lowest level and is equivalent to [CPU offloading](#cpu-offloading). But it can be made faster if you use streams without giving up inference speed.
```py
import torch
from diffusers import CogVideoXPipeline
from diffusers.hooks import apply_group_offloading
from diffusers.utils import export_to_video
onload_device = torch.device("cuda")
offload_device = torch.device("cpu")
pipeline = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b", torch_dtype=torch.bfloat16)
# Use the enable_group_offload method for Diffusers model implementations
pipeline.transformer.enable_group_offload(onload_device=onload_device, offload_device=offload_device, offload_type="leaf_level")
pipeline.vae.enable_group_offload(onload_device=onload_device, offload_type="leaf_level")
# Use the apply_group_offloading method for other model components
apply_group_offloading(pipeline.text_encoder, onload_device=onload_device, offload_type="block_level", num_blocks_per_group=2)
prompt = (
"A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. "
"The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other "
"pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, "
"casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. "
"The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical "
"atmosphere of this unique musical performance."
)
video = pipeline(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
export_to_video(video, "output.mp4", fps=8)
```
### CUDA stream
The `use_stream` parameter can be activated for CUDA devices that support asynchronous data transfer streams to reduce overall execution time compared to [CPU offloading](#cpu-offloading). It overlaps data transfer and computation by using layer prefetching. The next layer to be executed is loaded onto the GPU while the current layer is still being executed. It can increase CPU memory significantly so ensure you have 2x the amount of memory as the model size.
Set `record_stream=True` for more of a speedup at the cost of slightly increased memory usage. Refer to the [torch.Tensor.record_stream](https://pytorch.org/docs/stable/generated/torch.Tensor.record_stream.html) docs to learn more.
> [!TIP]
> When `use_stream=True` on VAEs with tiling enabled, make sure to do a dummy forward pass (possible with dummy inputs as well) before inference to avoid device mismatch errors. This may not work on all implementations, so feel free to open an issue if you encounter any problems.
If you're using `block_level` group offloading with `use_stream` enabled, the `num_blocks_per_group` parameter should be set to `1`, otherwise a warning will be raised.
```py
pipeline.transformer.enable_group_offload(onload_device=onload_device, offload_device=offload_device, offload_type="leaf_level", use_stream=True, record_stream=True)
```
The `low_cpu_mem_usage` parameter can be set to `True` to reduce CPU memory usage when using streams during group offloading. It is best for `leaf_level` offloading and when CPU memory is bottlenecked. Memory is saved by creating pinned tensors on the fly instead of pre-pinning them. However, this may increase overall execution time.
## Layerwise casting
Layerwise casting stores weights in a smaller data format (for example, `torch.float8_e4m3fn` and `torch.float8_e5m2`) to use less memory and upcasts those weights to a higher precision like `torch.float16` or `torch.bfloat16` for computation. Certain layers (normalization and modulation related weights) are skipped because storing them in fp8 can degrade generation quality.
> [!WARNING]
> Layerwise casting may not work with all models if the forward implementation contains internal typecasting of weights. The current implementation of layerwise casting assumes the forward pass is independent of the weight precision and the input datatypes are always specified in `compute_dtype` (see [here](https://github.com/huggingface/transformers/blob/7f5077e53682ca855afc826162b204ebf809f1f9/src/transformers/models/t5/modeling_t5.py#L294-L299) for an incompatible implementation).
>
> Layerwise casting may also fail on custom modeling implementations with [PEFT](https://huggingface.co/docs/peft/index) layers. There are some checks available but they are not extensively tested or guaranteed to work in all cases.
Call [`~ModelMixin.enable_layerwise_casting`] to set the storage and computation datatypes.
```py
import torch
from diffusers import CogVideoXPipeline, CogVideoXTransformer3DModel
from diffusers.utils import export_to_video
transformer = CogVideoXTransformer3DModel.from_pretrained(
"THUDM/CogVideoX-5b",
subfolder="transformer",
torch_dtype=torch.bfloat16
)
transformer.enable_layerwise_casting(storage_dtype=torch.float8_e4m3fn, compute_dtype=torch.bfloat16)
pipeline = CogVideoXPipeline.from_pretrained("THUDM/CogVideoX-5b",
transformer=transformer,
torch_dtype=torch.bfloat16
).to("cuda")
prompt = (
"A panda, dressed in a small, red jacket and a tiny hat, sits on a wooden stool in a serene bamboo forest. "
"The panda's fluffy paws strum a miniature acoustic guitar, producing soft, melodic tunes. Nearby, a few other "
"pandas gather, watching curiously and some clapping in rhythm. Sunlight filters through the tall bamboo, "
"casting a gentle glow on the scene. The panda's face is expressive, showing concentration and joy as it plays. "
"The background includes a small, flowing stream and vibrant green foliage, enhancing the peaceful and magical "
"atmosphere of this unique musical performance."
)
video = pipeline(prompt=prompt, guidance_scale=6, num_inference_steps=50).frames[0]
print(f"Max memory reserved: {torch.cuda.max_memory_allocated() / 1024**3:.2f} GB")
export_to_video(video, "output.mp4", fps=8)
```
The [`~hooks.apply_layerwise_casting`] method can also be used if you need more control and flexibility. It can be partially applied to model layers by calling it on specific internal modules. Use the `skip_modules_pattern` or `skip_modules_classes` parameters to specify modules to avoid, such as the normalization and modulation layers.
To use sliced VAE, call [`~StableDiffusionPipeline.enable_vae_slicing`] on your pipeline before inference:
```python
import torch
from diffusers import CogVideoXTransformer3DModel
from diffusers.hooks import apply_layerwise_casting
from diffusers import StableDiffusionPipeline
transformer = CogVideoXTransformer3DModel.from_pretrained(
"THUDM/CogVideoX-5b",
subfolder="transformer",
torch_dtype=torch.bfloat16
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
)
pipe = pipe.to("cuda")
# skip the normalization layer
apply_layerwise_casting(
transformer,
storage_dtype=torch.float8_e4m3fn,
compute_dtype=torch.bfloat16,
skip_modules_classes=["norm"],
non_blocking=True,
)
prompt = "a photo of an astronaut riding a horse on mars"
pipe.enable_vae_slicing()
#pipe.enable_xformers_memory_efficient_attention()
images = pipe([prompt] * 32).images
```
## torch.channels_last
You may see a small performance boost in VAE decoding on multi-image batches, and there should be no performance impact on single-image batches.
[torch.channels_last](https://pytorch.org/tutorials/intermediate/memory_format_tutorial.html) flips how tensors are stored from `(batch size, channels, height, width)` to `(batch size, heigh, width, channels)`. This aligns the tensors with how the hardware sequentially accesses the tensors stored in memory and avoids skipping around in memory to access the pixel values.
## Tiled VAE
Not all operators currently support the channels-last format and may result in worst performance, but it is still worth trying.
Tiled VAE processing also enables working with large images on limited VRAM (for example, generating 4k images on 8GB of VRAM) by splitting the image into overlapping tiles, decoding the tiles, and then blending the outputs together to compose the final image. You should also used tiled VAE with [`~ModelMixin.enable_xformers_memory_efficient_attention`] to reduce memory use further if you have xFormers installed.
```py
print(pipeline.unet.conv_out.state_dict()["weight"].stride()) # (2880, 9, 3, 1)
pipeline.unet.to(memory_format=torch.channels_last) # in-place operation
To use tiled VAE processing, call [`~StableDiffusionPipeline.enable_vae_tiling`] on your pipeline before inference:
```python
import torch
from diffusers import StableDiffusionPipeline, UniPCMultistepScheduler
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
)
pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
pipe = pipe.to("cuda")
prompt = "a beautiful landscape photograph"
pipe.enable_vae_tiling()
#pipe.enable_xformers_memory_efficient_attention()
image = pipe([prompt], width=3840, height=2224, num_inference_steps=20).images[0]
```
The output image has some tile-to-tile tone variation because the tiles are decoded separately, but you shouldn't see any sharp and obvious seams between the tiles. Tiling is turned off for images that are 512x512 or smaller.
## CPU offloading
Offloading the weights to the CPU and only loading them on the GPU when performing the forward pass can also save memory. Often, this technique can reduce memory consumption to less than 3GB.
To perform CPU offloading, call [`~StableDiffusionPipeline.enable_sequential_cpu_offload`]:
```Python
import torch
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
)
prompt = "a photo of an astronaut riding a horse on mars"
pipe.enable_sequential_cpu_offload()
image = pipe(prompt).images[0]
```
CPU offloading works on submodules rather than whole models. This is the best way to minimize memory consumption, but inference is much slower due to the iterative nature of the diffusion process. The UNet component of the pipeline runs several times (as many as `num_inference_steps`); each time, the different UNet submodules are sequentially onloaded and offloaded as needed, resulting in a large number of memory transfers.
<Tip>
Consider using [model offloading](#model-offloading) if you want to optimize for speed because it is much faster. The tradeoff is your memory savings won't be as large.
</Tip>
<Tip warning={true}>
When using [`~StableDiffusionPipeline.enable_sequential_cpu_offload`], don't move the pipeline to CUDA beforehand or else the gain in memory consumption will only be minimal (see this [issue](https://github.com/huggingface/diffusers/issues/1934) for more information).
[`~StableDiffusionPipeline.enable_sequential_cpu_offload`] is a stateful operation that installs hooks on the models.
</Tip>
## Model offloading
<Tip>
Model offloading requires 🤗 Accelerate version 0.17.0 or higher.
</Tip>
[Sequential CPU offloading](#cpu-offloading) preserves a lot of memory but it makes inference slower because submodules are moved to GPU as needed, and they're immediately returned to the CPU when a new module runs.
Full-model offloading is an alternative that moves whole models to the GPU, instead of handling each model's constituent *submodules*. There is a negligible impact on inference time (compared with moving the pipeline to `cuda`), and it still provides some memory savings.
During model offloading, only one of the main components of the pipeline (typically the text encoder, UNet and VAE)
is placed on the GPU while the others wait on the CPU. Components like the UNet that run for multiple iterations stay on the GPU until they're no longer needed.
Enable model offloading by calling [`~StableDiffusionPipeline.enable_model_cpu_offload`] on the pipeline:
```Python
import torch
from diffusers import StableDiffusionPipeline
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
)
prompt = "a photo of an astronaut riding a horse on mars"
pipe.enable_model_cpu_offload()
image = pipe(prompt).images[0]
```
<Tip warning={true}>
In order to properly offload models after they're called, it is required to run the entire pipeline and models are called in the pipeline's expected order. Exercise caution if models are reused outside the context of the pipeline after hooks have been installed. See [Removing Hooks](https://huggingface.co/docs/accelerate/en/package_reference/big_modeling#accelerate.hooks.remove_hook_from_module) for more information.
[`~StableDiffusionPipeline.enable_model_cpu_offload`] is a stateful operation that installs hooks on the models and state on the pipeline.
</Tip>
## Channels-last memory format
The channels-last memory format is an alternative way of ordering NCHW tensors in memory to preserve dimension ordering. Channels-last tensors are ordered in such a way that the channels become the densest dimension (storing images pixel-per-pixel). Since not all operators currently support the channels-last format, it may result in worst performance but you should still try and see if it works for your model.
For example, to set the pipeline's UNet to use the channels-last format:
```python
print(pipe.unet.conv_out.state_dict()["weight"].stride()) # (2880, 9, 3, 1)
pipe.unet.to(memory_format=torch.channels_last) # in-place operation
print(
pipeline.unet.conv_out.state_dict()["weight"].stride()
pipe.unet.conv_out.state_dict()["weight"].stride()
) # (2880, 1, 960, 320) having a stride of 1 for the 2nd dimension proves that it works
```
## torch.jit.trace
## Tracing
[torch.jit.trace](https://pytorch.org/docs/stable/generated/torch.jit.trace.html) records the operations a model performs on a sample input and creates a new, optimized representation of the model based on the recorded execution path. During tracing, the model is optimized to reduce overhead from Python and dynamic control flows and operations are fused together for more efficiency. The returned executable or [ScriptFunction](https://pytorch.org/docs/stable/generated/torch.jit.ScriptFunction.html) can be compiled.
Tracing runs an example input tensor through the model and captures the operations that are performed on it as that input makes its way through the model's layers. The executable or `ScriptFunction` that is returned is optimized with just-in-time compilation.
```py
To trace a UNet:
```python
import time
import torch
from diffusers import StableDiffusionPipeline
@@ -389,7 +191,8 @@ torch.set_grad_enabled(False)
n_experiments = 2
unet_runs_per_experiment = 50
# load sample inputs
# load inputs
def generate_inputs():
sample = torch.randn((2, 4, 64, 64), device="cuda", dtype=torch.float16)
timestep = torch.rand(1, device="cuda", dtype=torch.float16) * 999
@@ -397,12 +200,12 @@ def generate_inputs():
return sample, timestep, encoder_hidden_states
pipeline = StableDiffusionPipeline.from_pretrained(
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
).to("cuda")
unet = pipeline.unet
unet = pipe.unet
unet.eval()
unet.to(memory_format=torch.channels_last) # use channels_last memory format
unet.forward = functools.partial(unet.forward, return_dict=False) # set return_dict=False as default
@@ -419,12 +222,14 @@ unet_traced = torch.jit.trace(unet, inputs)
unet_traced.eval()
print("done tracing")
# warmup and optimize graph
for _ in range(5):
with torch.inference_mode():
inputs = generate_inputs()
orig_output = unet_traced(*inputs)
# benchmarking
with torch.inference_mode():
for _ in range(n_experiments):
@@ -446,18 +251,20 @@ with torch.inference_mode():
unet_traced.save("unet_traced.pt")
```
Replace the pipeline's UNet with the traced version.
Replace the `unet` attribute of the pipeline with the traced model:
```py
import torch
```python
from diffusers import StableDiffusionPipeline
import torch
from dataclasses import dataclass
@dataclass
class UNet2DConditionOutput:
sample: torch.Tensor
pipeline = StableDiffusionPipeline.from_pretrained(
pipe = StableDiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
@@ -466,7 +273,8 @@ pipeline = StableDiffusionPipeline.from_pretrained(
# use jitted unet
unet_traced = torch.jit.load("unet_traced.pt")
# del pipeline.unet
# del pipe.unet
class TracedUNet(torch.nn.Module):
def __init__(self):
super().__init__()
@@ -477,7 +285,8 @@ class TracedUNet(torch.nn.Module):
sample = unet_traced(latent_model_input, t, encoder_hidden_states)[0]
return UNet2DConditionOutput(sample=sample)
pipeline.unet = TracedUNet()
pipe.unet = TracedUNet()
with torch.inference_mode():
image = pipe([prompt] * 1, num_inference_steps=50).images[0]
@@ -485,31 +294,39 @@ with torch.inference_mode():
## Memory-efficient attention
> [!TIP]
> Memory-efficient attention optimizes for memory usage *and* [inference speed](./fp16#scaled-dot-product-attention!
Recent work on optimizing bandwidth in the attention block has generated huge speed-ups and reductions in GPU memory usage. The most recent type of memory-efficient attention is [Flash Attention](https://arxiv.org/abs/2205.14135) (you can check out the original code at [HazyResearch/flash-attention](https://github.com/HazyResearch/flash-attention)).
The Transformers attention mechanism is memory-intensive, especially for long sequences, so you can try using different and more memory-efficient attention types.
<Tip>
By default, if PyTorch >= 2.0 is installed, [scaled dot-product attention (SDPA)](https://pytorch.org/docs/stable/generated/torch.nn.functional.scaled_dot_product_attention.html) is used. You don't need to make any additional changes to your code.
If you have PyTorch >= 2.0 installed, you should not expect a speed-up for inference when enabling `xformers`.
SDPA supports [FlashAttention](https://github.com/Dao-AILab/flash-attention) and [xFormers](https://github.com/facebookresearch/xformers) as well as a native C++ PyTorch implementation. It automatically selects the most optimal implementation based on your input.
</Tip>
You can explicitly use xFormers with the [`~ModelMixin.enable_xformers_memory_efficient_attention`] method.
To use Flash Attention, install the following:
```py
# pip install xformers
- PyTorch > 1.12
- CUDA available
- [xFormers](xformers)
Then call [`~ModelMixin.enable_xformers_memory_efficient_attention`] on the pipeline:
```python
from diffusers import DiffusionPipeline
import torch
from diffusers import StableDiffusionXLPipeline
pipeline = StableDiffusionXLPipeline.from_pretrained(
"stabilityai/stable-diffusion-xl-base-1.0",
pipe = DiffusionPipeline.from_pretrained(
"stable-diffusion-v1-5/stable-diffusion-v1-5",
torch_dtype=torch.float16,
use_safetensors=True,
).to("cuda")
pipeline.enable_xformers_memory_efficient_attention()
pipe.enable_xformers_memory_efficient_attention()
with torch.inference_mode():
sample = pipe("a small cat")
# optional: You can disable it via
# pipe.disable_xformers_memory_efficient_attention()
```
Call [`~ModelMixin.disable_xformers_memory_efficient_attention`] to disable it.
```py
pipeline.disable_xformers_memory_efficient_attention()
```
The iteration speed when using `xformers` should match the iteration speed of PyTorch 2.0 as described [here](torch2.0).

View File

@@ -12,9 +12,6 @@ specific language governing permissions and limitations under the License.
# Metal Performance Shaders (MPS)
> [!TIP]
> Pipelines with a <img alt="MPS" src="https://img.shields.io/badge/MPS-000000?style=flat&logo=apple&logoColor=white%22"> badge indicate a model can take advantage of the MPS backend on Apple silicon devices for faster inference. Feel free to open a [Pull Request](https://github.com/huggingface/diffusers/compare) to add this badge to pipelines that are missing it.
🤗 Diffusers is compatible with Apple silicon (M1/M2 chips) using the PyTorch [`mps`](https://pytorch.org/docs/stable/notes/mps.html) device, which uses the Metal framework to leverage the GPU on MacOS devices. You'll need to have:
- macOS computer with Apple silicon (M1/M2) hardware
@@ -40,7 +37,7 @@ image
<Tip warning={true}>
The PyTorch [mps](https://pytorch.org/docs/stable/notes/mps.html) backend does not support NDArray sizes greater than `2**32`. Please open an [Issue](https://github.com/huggingface/diffusers/issues/new/choose) if you encounter this problem so we can investigate.
Generating multiple prompts in a batch can [crash](https://github.com/huggingface/diffusers/issues/363) or fail to work reliably. We believe this is related to the [`mps`](https://github.com/pytorch/pytorch/issues/84039) backend in PyTorch. While this is being investigated, you should iterate instead of batching.
</Tip>
@@ -62,10 +59,6 @@ If you're using **PyTorch 1.13**, you need to "prime" the pipeline with an addit
## Troubleshoot
This section lists some common issues with using the `mps` backend and how to solve them.
### Attention slicing
M1/M2 performance is very sensitive to memory pressure. When this occurs, the system automatically swaps if it needs to which significantly degrades performance.
To prevent this from happening, we recommend *attention slicing* to reduce memory pressure during inference and prevent swapping. This is especially relevant if your computer has less than 64GB of system RAM, or if you generate images at non-standard resolutions larger than 512×512 pixels. Call the [`~DiffusionPipeline.enable_attention_slicing`] function on your pipeline:
@@ -79,7 +72,3 @@ pipeline.enable_attention_slicing()
```
Attention slicing performs the costly attention operation in multiple steps instead of all at once. It usually improves performance by ~20% in computers without universal memory, but we've observed *better performance* in most Apple silicon computers unless you have 64GB of RAM or more.
### Batch inference
Generating multiple prompts in a batch can crash or fail to work reliably. If this is the case, try iterating instead of batching.

View File

@@ -1,497 +0,0 @@
# ParaAttention
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-performance.png">
</div>
<div class="flex justify-center">
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/hunyuan-video-performance.png">
</div>
Large image and video generation models, such as [FLUX.1-dev](https://huggingface.co/black-forest-labs/FLUX.1-dev) and [HunyuanVideo](https://huggingface.co/tencent/HunyuanVideo), can be an inference challenge for real-time applications and deployment because of their size.
[ParaAttention](https://github.com/chengzeyi/ParaAttention) is a library that implements **context parallelism** and **first block cache**, and can be combined with other techniques (torch.compile, fp8 dynamic quantization), to accelerate inference.
This guide will show you how to apply ParaAttention to FLUX.1-dev and HunyuanVideo on NVIDIA L20 GPUs.
No optimizations are applied for our baseline benchmark, except for HunyuanVideo to avoid out-of-memory errors.
Our baseline benchmark shows that FLUX.1-dev is able to generate a 1024x1024 resolution image in 28 steps in 26.36 seconds, and HunyuanVideo is able to generate 129 frames at 720p resolution in 30 steps in 3675.71 seconds.
> [!TIP]
> For even faster inference with context parallelism, try using NVIDIA A100 or H100 GPUs (if available) with NVLink support, especially when there is a large number of GPUs.
## First Block Cache
Caching the output of the transformers blocks in the model and reusing them in the next inference steps reduces the computation cost and makes inference faster.
However, it is hard to decide when to reuse the cache to ensure quality generated images or videos. ParaAttention directly uses the **residual difference of the first transformer block output** to approximate the difference among model outputs. When the difference is small enough, the residual difference of previous inference steps is reused. In other words, the denoising step is skipped.
This achieves a 2x speedup on FLUX.1-dev and HunyuanVideo inference with very good quality.
<figure>
<img src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/ada-cache.png" alt="Cache in Diffusion Transformer" />
<figcaption>How AdaCache works, First Block Cache is a variant of it</figcaption>
</figure>
<hfoptions id="first-block-cache">
<hfoption id="FLUX-1.dev">
To apply first block cache on FLUX.1-dev, call `apply_cache_on_pipe` as shown below. 0.08 is the default residual difference value for FLUX models.
```python
import time
import torch
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained(
"black-forest-labs/FLUX.1-dev",
torch_dtype=torch.bfloat16,
).to("cuda")
from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
apply_cache_on_pipe(pipe, residual_diff_threshold=0.08)
# Enable memory savings
# pipe.enable_model_cpu_offload()
# pipe.enable_sequential_cpu_offload()
begin = time.time()
image = pipe(
"A cat holding a sign that says hello world",
num_inference_steps=28,
).images[0]
end = time.time()
print(f"Time: {end - begin:.2f}s")
print("Saving image to flux.png")
image.save("flux.png")
```
| Optimizations | Original | FBCache rdt=0.06 | FBCache rdt=0.08 | FBCache rdt=0.10 | FBCache rdt=0.12 |
| - | - | - | - | - | - |
| Preview | ![Original](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-original.png) | ![FBCache rdt=0.06](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-fbc-0.06.png) | ![FBCache rdt=0.08](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-fbc-0.08.png) | ![FBCache rdt=0.10](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-fbc-0.10.png) | ![FBCache rdt=0.12](https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/flux-fbc-0.12.png) |
| Wall Time (s) | 26.36 | 21.83 | 17.01 | 16.00 | 13.78 |
First Block Cache reduced the inference speed to 17.01 seconds compared to the baseline, or 1.55x faster, while maintaining nearly zero quality loss.
</hfoption>
<hfoption id="HunyuanVideo">
To apply First Block Cache on HunyuanVideo, `apply_cache_on_pipe` as shown below. 0.06 is the default residual difference value for HunyuanVideo models.
```python
import time
import torch
from diffusers import HunyuanVideoPipeline, HunyuanVideoTransformer3DModel
from diffusers.utils import export_to_video
model_id = "tencent/HunyuanVideo"
transformer = HunyuanVideoTransformer3DModel.from_pretrained(
model_id,
subfolder="transformer",
torch_dtype=torch.bfloat16,
revision="refs/pr/18",
)
pipe = HunyuanVideoPipeline.from_pretrained(
model_id,
transformer=transformer,
torch_dtype=torch.float16,
revision="refs/pr/18",
).to("cuda")
from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
apply_cache_on_pipe(pipe, residual_diff_threshold=0.6)
pipe.vae.enable_tiling()
begin = time.time()
output = pipe(
prompt="A cat walks on the grass, realistic",
height=720,
width=1280,
num_frames=129,
num_inference_steps=30,
).frames[0]
end = time.time()
print(f"Time: {end - begin:.2f}s")
print("Saving video to hunyuan_video.mp4")
export_to_video(output, "hunyuan_video.mp4", fps=15)
```
<video controls>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/hunyuan-video-original.mp4" type="video/mp4">
Your browser does not support the video tag.
</video>
<small> HunyuanVideo without FBCache </small>
<video controls>
<source src="https://huggingface.co/datasets/huggingface/documentation-images/resolve/main/diffusers/para-attn/hunyuan-video-fbc.mp4" type="video/mp4">
Your browser does not support the video tag.
</video>
<small> HunyuanVideo with FBCache </small>
First Block Cache reduced the inference speed to 2271.06 seconds compared to the baseline, or 1.62x faster, while maintaining nearly zero quality loss.
</hfoption>
</hfoptions>
## fp8 quantization
fp8 with dynamic quantization further speeds up inference and reduces memory usage. Both the activations and weights must be quantized in order to use the 8-bit [NVIDIA Tensor Cores](https://www.nvidia.com/en-us/data-center/tensor-cores/).
Use `float8_weight_only` and `float8_dynamic_activation_float8_weight` to quantize the text encoder and transformer model.
The default quantization method is per tensor quantization, but if your GPU supports row-wise quantization, you can also try it for better accuracy.
Install [torchao](https://github.com/pytorch/ao/tree/main) with the command below.
```bash
pip3 install -U torch torchao
```
[torch.compile](https://pytorch.org/tutorials/intermediate/torch_compile_tutorial.html) with `mode="max-autotune-no-cudagraphs"` or `mode="max-autotune"` selects the best kernel for performance. Compilation can take a long time if it's the first time the model is called, but it is worth it once the model has been compiled.
This example only quantizes the transformer model, but you can also quantize the text encoder to reduce memory usage even more.
> [!TIP]
> Dynamic quantization can significantly change the distribution of the model output, so you need to change the `residual_diff_threshold` to a larger value for it to take effect.
<hfoptions id="fp8-quantization">
<hfoption id="FLUX-1.dev">
```python
import time
import torch
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained(
"black-forest-labs/FLUX.1-dev",
torch_dtype=torch.bfloat16,
).to("cuda")
from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
apply_cache_on_pipe(
pipe,
residual_diff_threshold=0.12, # Use a larger value to make the cache take effect
)
from torchao.quantization import quantize_, float8_dynamic_activation_float8_weight, float8_weight_only
quantize_(pipe.text_encoder, float8_weight_only())
quantize_(pipe.transformer, float8_dynamic_activation_float8_weight())
pipe.transformer = torch.compile(
pipe.transformer, mode="max-autotune-no-cudagraphs",
)
# Enable memory savings
# pipe.enable_model_cpu_offload()
# pipe.enable_sequential_cpu_offload()
for i in range(2):
begin = time.time()
image = pipe(
"A cat holding a sign that says hello world",
num_inference_steps=28,
).images[0]
end = time.time()
if i == 0:
print(f"Warm up time: {end - begin:.2f}s")
else:
print(f"Time: {end - begin:.2f}s")
print("Saving image to flux.png")
image.save("flux.png")
```
fp8 dynamic quantization and torch.compile reduced the inference speed to 7.56 seconds compared to the baseline, or 3.48x faster.
</hfoption>
<hfoption id="HunyuanVideo">
```python
import time
import torch
from diffusers import HunyuanVideoPipeline, HunyuanVideoTransformer3DModel
from diffusers.utils import export_to_video
model_id = "tencent/HunyuanVideo"
transformer = HunyuanVideoTransformer3DModel.from_pretrained(
model_id,
subfolder="transformer",
torch_dtype=torch.bfloat16,
revision="refs/pr/18",
)
pipe = HunyuanVideoPipeline.from_pretrained(
model_id,
transformer=transformer,
torch_dtype=torch.float16,
revision="refs/pr/18",
).to("cuda")
from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
apply_cache_on_pipe(pipe)
from torchao.quantization import quantize_, float8_dynamic_activation_float8_weight, float8_weight_only
quantize_(pipe.text_encoder, float8_weight_only())
quantize_(pipe.transformer, float8_dynamic_activation_float8_weight())
pipe.transformer = torch.compile(
pipe.transformer, mode="max-autotune-no-cudagraphs",
)
# Enable memory savings
pipe.vae.enable_tiling()
# pipe.enable_model_cpu_offload()
# pipe.enable_sequential_cpu_offload()
for i in range(2):
begin = time.time()
output = pipe(
prompt="A cat walks on the grass, realistic",
height=720,
width=1280,
num_frames=129,
num_inference_steps=1 if i == 0 else 30,
).frames[0]
end = time.time()
if i == 0:
print(f"Warm up time: {end - begin:.2f}s")
else:
print(f"Time: {end - begin:.2f}s")
print("Saving video to hunyuan_video.mp4")
export_to_video(output, "hunyuan_video.mp4", fps=15)
```
A NVIDIA L20 GPU only has 48GB memory and could face out-of-memory (OOM) errors after compilation and if `enable_model_cpu_offload` isn't called because HunyuanVideo has very large activation tensors when running with high resolution and large number of frames. For GPUs with less than 80GB of memory, you can try reducing the resolution and number of frames to avoid OOM errors.
Large video generation models are usually bottlenecked by the attention computations rather than the fully connected layers. These models don't significantly benefit from quantization and torch.compile.
</hfoption>
</hfoptions>
## Context Parallelism
Context Parallelism parallelizes inference and scales with multiple GPUs. The ParaAttention compositional design allows you to combine Context Parallelism with First Block Cache and dynamic quantization.
> [!TIP]
> Refer to the [ParaAttention](https://github.com/chengzeyi/ParaAttention/tree/main) repository for detailed instructions and examples of how to scale inference with multiple GPUs.
If the inference process needs to be persistent and serviceable, it is suggested to use [torch.multiprocessing](https://pytorch.org/docs/stable/multiprocessing.html) to write your own inference processor. This can eliminate the overhead of launching the process and loading and recompiling the model.
<hfoptions id="context-parallelism">
<hfoption id="FLUX-1.dev">
The code sample below combines First Block Cache, fp8 dynamic quantization, torch.compile, and Context Parallelism for the fastest inference speed.
```python
import time
import torch
import torch.distributed as dist
from diffusers import FluxPipeline
dist.init_process_group()
torch.cuda.set_device(dist.get_rank())
pipe = FluxPipeline.from_pretrained(
"black-forest-labs/FLUX.1-dev",
torch_dtype=torch.bfloat16,
).to("cuda")
from para_attn.context_parallel import init_context_parallel_mesh
from para_attn.context_parallel.diffusers_adapters import parallelize_pipe
from para_attn.parallel_vae.diffusers_adapters import parallelize_vae
mesh = init_context_parallel_mesh(
pipe.device.type,
max_ring_dim_size=2,
)
parallelize_pipe(
pipe,
mesh=mesh,
)
parallelize_vae(pipe.vae, mesh=mesh._flatten())
from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
apply_cache_on_pipe(
pipe,
residual_diff_threshold=0.12, # Use a larger value to make the cache take effect
)
from torchao.quantization import quantize_, float8_dynamic_activation_float8_weight, float8_weight_only
quantize_(pipe.text_encoder, float8_weight_only())
quantize_(pipe.transformer, float8_dynamic_activation_float8_weight())
torch._inductor.config.reorder_for_compute_comm_overlap = True
pipe.transformer = torch.compile(
pipe.transformer, mode="max-autotune-no-cudagraphs",
)
# Enable memory savings
# pipe.enable_model_cpu_offload(gpu_id=dist.get_rank())
# pipe.enable_sequential_cpu_offload(gpu_id=dist.get_rank())
for i in range(2):
begin = time.time()
image = pipe(
"A cat holding a sign that says hello world",
num_inference_steps=28,
output_type="pil" if dist.get_rank() == 0 else "pt",
).images[0]
end = time.time()
if dist.get_rank() == 0:
if i == 0:
print(f"Warm up time: {end - begin:.2f}s")
else:
print(f"Time: {end - begin:.2f}s")
if dist.get_rank() == 0:
print("Saving image to flux.png")
image.save("flux.png")
dist.destroy_process_group()
```
Save to `run_flux.py` and launch it with [torchrun](https://pytorch.org/docs/stable/elastic/run.html).
```bash
# Use --nproc_per_node to specify the number of GPUs
torchrun --nproc_per_node=2 run_flux.py
```
Inference speed is reduced to 8.20 seconds compared to the baseline, or 3.21x faster, with 2 NVIDIA L20 GPUs. On 4 L20s, inference speed is 3.90 seconds, or 6.75x faster.
</hfoption>
<hfoption id="HunyuanVideo">
The code sample below combines First Block Cache and Context Parallelism for the fastest inference speed.
```python
import time
import torch
import torch.distributed as dist
from diffusers import HunyuanVideoPipeline, HunyuanVideoTransformer3DModel
from diffusers.utils import export_to_video
dist.init_process_group()
torch.cuda.set_device(dist.get_rank())
model_id = "tencent/HunyuanVideo"
transformer = HunyuanVideoTransformer3DModel.from_pretrained(
model_id,
subfolder="transformer",
torch_dtype=torch.bfloat16,
revision="refs/pr/18",
)
pipe = HunyuanVideoPipeline.from_pretrained(
model_id,
transformer=transformer,
torch_dtype=torch.float16,
revision="refs/pr/18",
).to("cuda")
from para_attn.context_parallel import init_context_parallel_mesh
from para_attn.context_parallel.diffusers_adapters import parallelize_pipe
from para_attn.parallel_vae.diffusers_adapters import parallelize_vae
mesh = init_context_parallel_mesh(
pipe.device.type,
)
parallelize_pipe(
pipe,
mesh=mesh,
)
parallelize_vae(pipe.vae, mesh=mesh._flatten())
from para_attn.first_block_cache.diffusers_adapters import apply_cache_on_pipe
apply_cache_on_pipe(pipe)
# from torchao.quantization import quantize_, float8_dynamic_activation_float8_weight, float8_weight_only
#
# torch._inductor.config.reorder_for_compute_comm_overlap = True
#
# quantize_(pipe.text_encoder, float8_weight_only())
# quantize_(pipe.transformer, float8_dynamic_activation_float8_weight())
# pipe.transformer = torch.compile(
# pipe.transformer, mode="max-autotune-no-cudagraphs",
# )
# Enable memory savings
pipe.vae.enable_tiling()
# pipe.enable_model_cpu_offload(gpu_id=dist.get_rank())
# pipe.enable_sequential_cpu_offload(gpu_id=dist.get_rank())
for i in range(2):
begin = time.time()
output = pipe(
prompt="A cat walks on the grass, realistic",
height=720,
width=1280,
num_frames=129,
num_inference_steps=1 if i == 0 else 30,
output_type="pil" if dist.get_rank() == 0 else "pt",
).frames[0]
end = time.time()
if dist.get_rank() == 0:
if i == 0:
print(f"Warm up time: {end - begin:.2f}s")
else:
print(f"Time: {end - begin:.2f}s")
if dist.get_rank() == 0:
print("Saving video to hunyuan_video.mp4")
export_to_video(output, "hunyuan_video.mp4", fps=15)
dist.destroy_process_group()
```
Save to `run_hunyuan_video.py` and launch it with [torchrun](https://pytorch.org/docs/stable/elastic/run.html).
```bash
# Use --nproc_per_node to specify the number of GPUs
torchrun --nproc_per_node=8 run_hunyuan_video.py
```
Inference speed is reduced to 649.23 seconds compared to the baseline, or 5.66x faster, with 8 NVIDIA L20 GPUs.
</hfoption>
</hfoptions>
## Benchmarks
<hfoptions id="conclusion">
<hfoption id="FLUX-1.dev">
| GPU Type | Number of GPUs | Optimizations | Wall Time (s) | Speedup |
| - | - | - | - | - |
| NVIDIA L20 | 1 | Baseline | 26.36 | 1.00x |
| NVIDIA L20 | 1 | FBCache (rdt=0.08) | 17.01 | 1.55x |
| NVIDIA L20 | 1 | FP8 DQ | 13.40 | 1.96x |
| NVIDIA L20 | 1 | FBCache (rdt=0.12) + FP8 DQ | 7.56 | 3.48x |
| NVIDIA L20 | 2 | FBCache (rdt=0.12) + FP8 DQ + CP | 4.92 | 5.35x |
| NVIDIA L20 | 4 | FBCache (rdt=0.12) + FP8 DQ + CP | 3.90 | 6.75x |
</hfoption>
<hfoption id="HunyuanVideo">
| GPU Type | Number of GPUs | Optimizations | Wall Time (s) | Speedup |
| - | - | - | - | - |
| NVIDIA L20 | 1 | Baseline | 3675.71 | 1.00x |
| NVIDIA L20 | 1 | FBCache | 2271.06 | 1.62x |
| NVIDIA L20 | 2 | FBCache + CP | 1132.90 | 3.24x |
| NVIDIA L20 | 4 | FBCache + CP | 718.15 | 5.12x |
| NVIDIA L20 | 8 | FBCache + CP | 649.23 | 5.66x |
</hfoption>
</hfoptions>

View File

@@ -78,23 +78,6 @@ For more information and different options about `torch.compile`, refer to the [
> [!TIP]
> Learn more about other ways PyTorch 2.0 can help optimize your model in the [Accelerate inference of text-to-image diffusion models](../tutorials/fast_diffusion) tutorial.
### Regional compilation
Compiling the whole model usually has a big problem space for optimization. Models are often composed of multiple repeated blocks. [Regional compilation](https://pytorch.org/tutorials/recipes/regional_compilation.html) compiles the repeated block first (a transformer encoder block, for example), so that the Torch compiler would re-use its cached/optimized generated code for the other blocks, reducing (often massively) the cold start compilation time observed on the first inference call.
Enabling regional compilation might require simple yet intrusive changes to the
modeling code. However, 🤗 Accelerate provides a utility [`compile_regions()`](https://huggingface.co/docs/accelerate/main/en/usage_guides/compilation#how-to-use-regional-compilation) which automatically compiles
the repeated blocks of the provided `nn.Module` sequentially, and the rest of the model separately. This helps with reducing cold start time while keeping most (if not all) of the speedup you would get from full compilation.
```py
# Make sure you're on the latest `accelerate`: `pip install -U accelerate`.
from accelerate.utils import compile_regions
pipe.unet = compile_regions(pipe.unet, mode="reduce-overhead", fullgraph=True)
```
As you may have noticed `compile_regions()` takes the same arguments as `torch.compile()`, allowing flexibility.
## Benchmark
We conducted a comprehensive benchmark with PyTorch 2.0's efficient attention implementation and `torch.compile` across different GPUs and batch sizes for five of our most used pipelines. The code is benchmarked on 🤗 Diffusers v0.17.0.dev0 to optimize `torch.compile` usage (see [here](https://github.com/huggingface/diffusers/pull/3313) for more details).

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